Nano's Site

Mencari uang dollar di onbux

2011-03-13T01:01:00.000-08:00

Sejarah lahirnya Onbux dimulai pada tanggal 01-Dec-08 namun di kenal eksis pada tanggal 11-Jun-10, namun pun demikian PTC ini telah membuktikan bahwa dia salah satu PTC terbaik dengan total payment sekitar 2,5 juta dollar, Sungguh angka yang fantastis di usianya yang masih terbilang muda sejak PTC ini eksis.

Manfaatkan Peluang bisnis ini dengan mendaftar di PTC ONBUX

CARA CEPAT MENGISI PAYPAL GRATIS

2011-03-12T20:14:00.000-08:00

Bagaimana Cara cepat mengisi paypal dengan dollar secara gratis? Sangatlah mudah.. dan uang yang bisa anda hasilkan adalah sebesar $75. mungkin cara ini sedikit kocak dan konyol, tp silahkan dicoba aja kalo gak percaya…! sebelum melanjutkan langkah-langkah berikut, sebaiknya anda sudah memiliki rekening e-currency paypal dulu.. jika belum punya paypal, silahkan daftar di sini..!

Trik ini adalah murni dari pengalaman pribadi, dimana kita akan memanfaatkan situs jejaring sosial sebagai media promosi dan A.W.Survey sebagai lahannya. dan untuk jika anda ini mengetahui lebih dalam tentang A.W.Survey, silahkan pelajari di postingan sebelumnya.

Adapun langkah-langkahnya adalah sebagai berikut :

Daftar dulu di A.W.SURVEY dengan meng-klik di sini.
setelah anda sukses mendaftar dan mendapatkan $27, anda hanya membutuhkan $48 lagi untuk mencapai Payout. dan itu artinya anda tinggal mengajak teman anda sebanyak 39 orang untuk mencapainya.
Sy yakin anda memiliki teman dekat lebih dari 39 orang..! nah inilah yang kita manfaatkan!
Langkah selanjutnya adalah, meminta teman dekat anda tersebut untuk bergabung di A.W.SURVEY dengan menggunakan link referral anda. Untuk melihat link referal anda, silahkan klik di sini (di bawah tulisan Your Direct Referring Link is below)
Untuk memperlancar misi anda, katakan pada teman anda bahwa anda lagi butuh duit dan mereka bisa membantu hanya dengan mendaftar di A.W.Survey. dan jangan lupa yakinkan mereka kalo ente lagi terdesak.
Ucapkan terima kasih kepada mereka semua yang telah membantumu..
Nah, jika anda tekun, anda hanya butu waktu kurang dari 2 jam mengerjakan hal ini dan uang sebesar $75 siap untuk di tarik (PO) ke Paypal anda.
Silahkan menggunakan metode chating atau kirim pesan ke teman anda untuk melakukan hal di atas.
Selamat mencoba…!

Catatan : dengan perbandingan 1:3 antara yang sukses dan gagal.. Sy sarankan melakukan perlakuan kepada 100 teman dekat anda. dan jika sukses, anda telah mendapatkan uang cuma-cuma sebesar $75 (Rp 750.000 / kurs 10.000).

Google Adsense

2011-03-07T21:08:00.001-08:00

Google AdSense is a new, fast and easy way for small and medium sites to make money free online by displaying relevant, text-based, un-obtrusive ads from Google AdWords (Google's own advertising program) and receiving a share of the pay-per-click payment.
Because the ads are related to what your users are looking for on your site, the results can be much better than you'd earn from banner networks and many affiliate programs.
For now, AdSense is the best way to make money free online from informational sites even if there are no obvious related affiliate programs. But you don't need to disregard affiliate programs. You can combine both these ways to make money free online and double your income.
AdSense is easy to join, it doesn't cost you anything, all you have to do is paste a few lines of code into your pages, and Google does the work of finding the best ads for them from hundreds of thousands of AdWords advertisers. You can check the relevance of the ads by looking at the text ads on the right side of this page.

How to get started

Go to Google AdSense.
Fill in the application form and confirm an email that Google will send you. If you own several sites, you need apply only once.
Google evaluates your site and will follow-up with you via email within 2-3 days (usually within 24 hours). If you're accepted, you'll be able to log in to your AdSense account.
Log in to your account using the email address and password that you submitted with your application, and agree to the AdSense Terms and Conditions.
Paste the AdSense ad code into your Web pages. There are 10 ad layout choices: 728x90, 468x60, 125x125, 120x600, 160x600, 120x240, 300x250, 250x250, 336x280, and 180x150. In addition to text ads, you have an option of running contextually targeted image ads.
You can choose the color palette from a long list of available palettes or create your own. You can even rotate your ads through up to 4 palettes.
The AdSense ad code is unique for your account and is not site/page-specific. You can place it on any page or site you own.

How to get the most out of AdSense

There are three obvious ways to increase your income from Google AdSense...
1) Increase traffic
Create more keyword-focused pages. See Building: Site Content: Choosing Keywords for more information on finding the most profitable keywords.
2) Increase click-through rate
Use simple design with the AdSense ads displayed prominently. According to Google, ads in the skyscraper format works better (especially on the right side of the page). Focus to only one topic per page - that should make it easier for Google to serve up the more tightly contextual ads which means better click-through.
3) Increase the value of clicks
Of course, you can't do it directly. However, you can find some "expensive" keywords and create pages optimized for them, within your site's theme. These keywords are highly competitive and you'll unlikely get high ranking for them, but visitors will arrive from other "inexpensive" pages and click on Google's ads.
To estimate the relative value of a keyword, search for it on FindWhat. The search results page will show the cost for each listing.

Capacitor

2010-05-16T21:02:00.000-07:00

A capacitor or condenser is a passive electronic component consisting of a pair of conductors separated by a dielectric (insulator). When a potential difference (voltage) exists across the conductors, an electric field is present in the dielectric. This field stores energy and produces a mechanical force between the conductors. The effect is greatest when there is a narrow separation between large areas of conductor, hence capacitor conductors are often called plates.
An ideal capacitor is characterized by a single constant value, capacitance, which is measured in farads. This is the ratio of the electric charge on each conductor to the potential difference between them. In practice, the dielectric between the plates passes a small amount of leakage current. The conductors and leads introduce an equivalent series resistance and the dielectric has an electric field strength limit resulting in a breakdown voltage.
Capacitors are widely used in electronic circuits to block the flow of direct current while allowing alternating current to pass, to filter out interference, to smooth the output of power supplies, and for many other purposes. They are used in resonant circuits in radio frequency equipment to select particular frequencies from a signal with many frequencies.

History

Battery of four Leyden jars in Museum Boerhave, Leiden, the Netherlands.

In October 1745, Ewald Georg von Kleist of Pomerania in Germany found that charge could be stored by connecting a high voltage electrostatic generator by a wire to a volume of water in a hand-held glass jar.^[1] Von Kleist's hand and the water acted as conductors and the jar as a dielectric (although details of the mechanism were incorrectly identified at the time). Von Kleist found that after removing the generator, touching the wire resulted in a painful spark. In a letter describing the experiment, he said "I would not take a second shock for the kingdom of France."^[2] The following year, the Dutch physicist Pieter van Musschenbroek invented a similar capacitor, which was named the Leyden jar, after the University of Leiden where he worked.^[3] Daniel Gralath was the first to combine several jars in parallel into a "battery" to increase the charge storage capacity.^{[citation needed]}
Benjamin Franklin investigated the Leyden jar and "proved" that the charge was stored on the glass, not in the water as others had assumed. He also created the term "battery",^[4]^[5] (as in a battery of cannon), subsequently applied to clusters of electrochemical cells.^[6] Leyden jars were later to be made by coating the inside and outside of jars with metal foil, leaving a space at the mouth to prevent arcing between the foils.^{[citation needed]} The earliest unit of capacitance was the 'jar', equivalent to about 1 nanofarad.^{[citation needed]}
Leyden jars or more powerful devices employing flat glass plates alternating with foil conductors were used exclusively up until about 1900, when the invention of wireless (radio) created a demand for standard capacitors, and the steady move to higher frequencies required capacitors with lower inductance. A more compact construction began to be used of a flexible dielectric sheet such as oiled paper sandwiched between sheets of metal foil, rolled or folded into a small package.
Early capacitors were also known as condensers, a term that is still occasionally used today. The term was first used for this purpose by Alessandro Volta in 1782, with reference to the device's ability to store a higher density of electric charge than a normal isolated conductor.^{[citation needed]}

Theory of operation

Main article: Capacitance

Charge separation in a parallel-plate capacitor causes an internal electric field. A dielectric (orange) reduces the field and increases the capacitance.

A simple demonstration of a parallel-plate capacitor

A capacitor consists of two conductors separated by a non-conductive region.^[7] The non-conductive substance is called the dielectric medium, although this may also mean a vacuum or a semiconductor depletion region chemically identical to the conductors. A capacitor is assumed to be self-contained and isolated, with no net electric charge and no influence from an external electric field. The conductors thus contain equal and opposite charges on their facing surfaces,^[8] and the dielectric contains an electric field. The capacitor is a reasonably general model for electric fields within electric circuits.
An ideal capacitor is wholly characterized by a constant capacitance C, defined as the ratio of charge ±Q on each conductor to the voltage V between them:^[7]

Sometimes charge buildup affects the mechanics of the capacitor, causing the capacitance to vary. In this case, capacitance is defined in terms of incremental changes:

In SI units, a capacitance of one farad means that one coulomb of charge on each conductor causes a voltage of one volt across the device.^[9]

Energy storage

Work must be done by an external influence to move charge between the conductors in a capacitor. When the external influence is removed, the charge separation persists and energy is stored in the electric field. If charge is later allowed to return to its equilibrium position, the energy is released. The work done in establishing the electric field, and hence the amount of energy stored, is given by:^[10]

Current-voltage relation

The current i(t) through a component in an electric circuit is defined as the rate of flow of the charge q(t) that has passed through it. Physical charges cannot pass through the dielectric layer of a capacitor, but rather build up in equal and opposite quantities on the electrodes: as each electron accumulates on the negative plate, one leaves the positive plate. Thus the accumulated charge on the electrodes is equal to the integral of the current, as well as being proportional to the voltage (as discussed above). As with any antiderivative, a constant of integration is added to represent the initial voltage v (t₀). This is the integral form of the capacitor equation,^[11]

Taking the derivative of this, and multiplying by C, yields the derivative form,^[12]

The dual of the capacitor is the inductor, which stores energy in the magnetic field rather than the electric field. Its current-voltage relation is obtained by exchanging current and voltage in the capacitor equations and replacing C with the inductance L.

DC circuits

AC circuits

Impedance, the vector sum of reactance and resistance, describes the phase difference and the ratio of amplitudes between sinusoidally varying voltage and sinusoidally varying current at a given frequency. Fourier analysis allows any signal to be constructed from a spectrum of frequencies, whence the circuit's reaction to the various frequencies may be found. The reactance and impedance of a capacitor are respectively

where j is the imaginary unit and ω is the angular velocity of the sinusoidal signal. The - j phase indicates that the AC voltage V = Z I lags the AC current by 90°: the positive current phase corresponds to increasing voltage as the capacitor charges; zero current corresponds to instantaneous constant voltage, etc.
Note that impedance decreases with increasing capacitance and increasing frequency. This implies that a higher-frequency signal or a larger capacitor results in a lower voltage amplitude per current amplitude—an AC "short circuit" or AC coupling. Conversely, for very low frequencies, the reactance will be high, so that a capacitor is nearly an open circuit in AC analysis—those frequencies have been "filtered out".
Capacitors are different from resistors and inductors in that the impedance is inversely proportional to the defining characteristic, i.e. capacitance.

Parallel plate model

Dielectric is placed between two conducting plates, each of area A and with a separation of d.

The simplest capacitor consists of two parallel conductive plates separated by a dielectric with permittivity ε (such as air). The model may also be used to make qualitative predictions for other device geometries. The plates are considered to extend uniformly over an area A and a charge density ±ρ = ±Q/A exists on their surface. Assuming that the width of the plates is much greater than their separation d, the electric field near the centre of the device will be uniform with the magnitude E = ρ/ε. The voltage is defined as the line integral of the electric field between the plates

Solving this for C = Q/V reveals that capacitance increases with area and decreases with separation

The capacitance is therefore greatest in devices made from materials with a high permittivity.

Several capacitors in parallel.

Networks

Non-ideal behaviour

Capacitors deviate from the ideal capacitor equation in a number of ways. Some of these, such as leakage current and parasitic effects are linear, or can be assumed to be linear, and can be dealt with by adding virtual components to the equivalent circuit of the capacitor. The usual methods of network analysis can then be applied. In other cases, such as with breakdown voltage, the effect is non-linear and normal (i.e., linear) network analysis cannot be used, the effect must be dealt with separately. There is yet another group, which may be linear but invalidate the assumption in the analysis that capacitance is a constant. Such an example is temperature dependence.

Breakdown voltage

Main article: Breakdown voltage

Above a particular electric field, known as the dielectric strength E_ds, the dielectric in a capacitor becomes conductive. The voltage at which this occurs is called the breakdown voltage of the device, and is given by the product of the dielectric strength and the separation between the conductors,^[14]

V bd = E ds d

The maximum energy that can be stored safely in a capacitor is limited by the breakdown voltage. Due to the scaling of capacitance and breakdown voltage with dielectric thickness, all capacitors made with a particular dielectric have approximately equal maximum energy density, to the extent that the dielectric dominates their volume.^[15]
For air dielectric capacitors the breakdown field strength is of the order 10⁷ V/m and will be much less when other materials are used for the dielectric. The absolute breakdown voltage of most capacitors is nowhere near such a high number because of the very small distance between the plates. Typical ratings for capacitors used for general electronics applications range from a few volts to 100V or so. For high voltage applications physically much larger capacitors have to be used. In this field, there are a number of factors that can dramatically reduce the breakdown voltage below the value to be expected by considering the breakdown field strength of the dielectric alone. For one thing, the geometry of the capacitor conductive parts (plates and connecting wires) is important. In particular, sharp edges or points hugely increase the electric field strength at that point and can lead to a local breakdown. Once this starts to happen, the breakdown will quickly "track" through the dielectric till it reaches the opposite plate and cause a short circuit.^[16]
The usual breakdown route is that the field strength becomes large enough to pull electrons in the dielectric from their atoms thus causing conduction. Other scenarios are possible, such as impurities in the dielectric, and, if the dielectric is of a crystalline nature, imperfections in the crystal structure can result in an avalanche breakdown as seen in semi-conductor devices. Breakdown voltage is also affected by pressure, humidity and temperature.^[17]

Equivalent circuit

Two equivalent circuits of a real capacitor

An ideal capacitor only stores and releases electrical energy, without dissipating any. In reality, all capacitors have imperfections within the capacitor's material that create resistance. This is specified as the equivalent series resistance or ESR of a component. This adds a real component to the impedance:

As frequency approaches infinity, the capacitive impedance (or reactance) approaches zero and the ESR becomes significant. As the reactance becomes negligible, power dissipation approaches P_RMS. = V_RMS.² /R_ESR.
Similarly to ESR, the capacitor's leads add equivalent series inductance or ESL to the component. This is usually significant only at relatively high frequencies. As inductive reactance is positive and increases with frequency, above a certain frequency capacitance will be canceled by inductance. High frequency engineering involves accounting for the inductance of all connections and components.
If the conductors are separated by a material with a small conductivity rather than a perfect dielectric, then a small leakage current flows directly between them. The capacitor therefore has a finite parallel resistance,^[9] and slowly discharges over time (time may vary greatly depending on the capacitor material and quality).

Ripple current

Ripple current is the AC component of an applied source (often a switched-mode power supply) whose frequency may be constant or varying. Certain types of capacitors, such as electrolytic tantalum capacitors, usually have a rating for maximum ripple current (both in frequency and magnitude). This ripple current can cause damaging heat to be generated within the capacitor due to the current flow across resistive imperfections in the materials used within the capacitor, more commonly referred to as equivalent series resistance (ESR). For example electrolytic tantalum capacitors are limited by ripple current and generally have the highest ESR ratings in the capacitor family, while ceramic capacitors generally have no ripple current limitation and have some of the lowest ESR ratings.

Capacitance instability

The capacitance of certain capacitors decreases as the component ages. In ceramic capacitors, this is caused by degradation of the dielectric. The type of dielectric and the ambient operating and storage temperatures are the most significant aging factors, while the operating voltage has a smaller effect. The aging process may be reversed by heating the component above the Curie point. Aging is fastest near the beginning of life of the component, and the device stabilizes over time.^[18] Electrolytic capacitors age as the electrolyte evaporates. In contrast with ceramic capacitors, this occurs towards the end of life of the component.
Temperature dependence of capacitance is usually expressed in parts per million (ppm) per °C. It can usually be taken as a broadly linear function but can be noticeably non-linear at the temperature extremes. The temperature coefficient can be either positive or negative, sometimes even amongst different samples of the same type. In other words, the spread in the range of temperature coefficients can encompass zero. See the data sheet in the leakage current section above for an example.
Capacitors, especially older components, can absorb sound waves resulting in a microphonic effect. Vibration moves the plates, causing the capacitance to vary, in turn inducing AC current. Some dielectrics also generate piezoelectricity. The resulting interference is especially problematic in audio applications, potentially causing feedback or unintended recording. In the reverse microphonic effect, the varying electric field between the capacitor plates exerts a physical force, moving them as a speaker. This can generate audible sound, but drains energy and stresses the dielectric and the electrolyte, if any.

Capacitor types

Main article: Types of capacitor

Practical capacitors are available commercially in many different forms. The type of internal dielectric, the structure of the plates and the device packaging all strongly affect the characteristics of the capacitor, and its applications.
Values available range from very low (picofarad range; while arbitrarily low values are in principle possible, stray (parasitic) capacitance in any circuit is the limiting factor) to about 5 kF supercapacitors.
Above approximately 1 microfarad electrolytic capacitors are usually used because of their small size and low cost compared with other technologies, unless their relatively poor stability, life and polarised nature make them unsuitable. Very high capacity supercapacitors use a porous carbon-based electrode material.

Dielectric materials

Capacitor materials. From left: multilayer ceramic, ceramic disc, multilayer polyester film, tubular ceramic, polystyrene, metalized polyester film, aluminum electrolytic. Major scale divisions are in centimetres.

Most types of capacitor include a dielectric spacer, which increases their capacitance. These dielectrics are most often insulators. However, low capacitance devices are available with a vacuum between their plates, which allows extremely high voltage operation and low losses. Variable capacitors with their plates open to the atmosphere were commonly used in radio tuning circuits. Later designs use polymer foil dielectric between the moving and stationary plates, with no significant air space between them.
In order to maximise the charge that a capacitor can hold, the dialectric material needs to have as high a permittivity as possible, while also having as high a breakdown voltage as possible.
Several solid dielectrics are available, including paper, plastic, glass, mica and ceramic materials. Paper was used extensively in older devices and offers relatively high voltage performance. However, it is susceptible to water absorption, and has been largely replaced by plastic film capacitors. Plastics offer better stability and aging performance, which makes them useful in timer circuits, although they may be limited to low operating temperatures and frequencies. Ceramic capacitors are generally small, cheap and useful for high frequency applications, although their capacitance varies strongly with voltage and they age poorly. They are broadly categorized as class 1 dielectrics, which have predictable variation of capacitance with temperature or class 2 dielectrics, which can operate at higher voltage. Glass and mica capacitors are extremely reliable, stable and tolerant to high temperatures and voltages, but are too expensive for most mainstream applications. Electrolytic capacitors and supercapacitors are used to store small and larger amounts of energy, respectively, ceramic capacitors are often used in resonators, and parasitic capacitance occurs in circuits wherever the simple conductor-insulator-conductor structure is formed unintentionally by the configuration of the circuit layout.
Electrolytic capacitors use an aluminum or tantalum plate with an oxide dielectric layer. The second electrode is a liquid electrolyte, connected to the circuit by another foil plate. Electrolytic capacitors offer very high capacitance but suffer from poor tolerances, high instability, gradual loss of capacitance especially when subjected to heat, and high leakage current. The conductivity of the electrolyte drops at low temperatures, which increases equivalent series resistance. While widely used for power-supply conditioning, poor high-frequency characteristics make them unsuitable for many applications. Tantalum capacitors offer better frequency and temperature characteristics than aluminum, but higher dielectric absorption and leakage.^[19] OS-CON (or OC-CON) capacitors are a polymerized organic semiconductor solid-electrolyte type that offer longer life at higher cost than standard electrolytic capacitors.
Several other types of capacitor are available for specialist applications. Supercapacitors store large amounts of energy. Supercapacitors made from carbon aerogel, carbon nanotubes, or highly porous electrode materials offer extremely high capacitance (up to 5 kF as of 2010) and can be used in some applications instead of rechargeable batteries. Alternating current capacitors are specifically designed to work on line (mains) voltage AC power circuits. They are commonly used in electric motor circuits and are often designed to handle large currents, so they tend to be physically large. They are usually ruggedly packaged, often in metal cases that can be easily grounded/earthed. They also are designed with direct current breakdown voltages of at least five times the maximum AC voltage.

Structure

Capacitor packages: SMD ceramic at top left; SMD tantalum at bottom left; through-hole tantalum at top right; through-hole electrolytic at bottom right. Major scale divisions are cm.

The arrangement of plates and dielectric has many variations depending on the desired ratings of the capacitor. For small values of capacitance (microfarads and less), ceramic disks use metallic coatings, with wire leads bonded to the coating. Larger values can be made by multiple stacks of plates and disks. Larger value capacitors usually use a metal foil or metal film layer deposited on the surface of a dielectric film to make the plates, and a dielectric film of impregnated paper or plastic – these are rolled up to save space. To reduce the series resistance and inductance for long plates, the plates and dielectric are staggered so that connection is made at the common edge of the rolled-up plates, not at the ends of the foil or metalized film strips that comprise the plates.
The assembly is encased to prevent moisture entering the dielectric – early radio equipment used a cardboard tube sealed with wax. Modern paper or film dielectric capacitors are dipped in a hard thermoplastic. Large capacitors for high-voltage use may have the roll form compressed to fit into a rectangular metal case, with bolted terminals and bushings for connections. The dielectric in larger capacitors is often impregnated with a liquid to improve its properties.
Capacitors may have their connecting leads arranged in many configurations, for example axially or radially. "Axial" means that the leads are on a common axis, typically the axis of the capacitor's cylindrical body – the leads extend from opposite ends. Radial leads might more accurately be referred to as tandem; they are rarely actually aligned along radii of the body's circle, so the term is inexact, although universal. The leads (until bent) are usually in planes parallel to that of the flat body of the capacitor, and extend in the same direction; they are often parallel as manufactured.
Small, cheap discoidal ceramic capacitors have existed since the 1930s, and remain in widespread use. Since the 1980s, surface mount packages for capacitors have been widely used. These packages are extremely small and lack connecting leads, allowing them to be soldered directly onto the surface of printed circuit boards. Surface mount components avoid undesirable high-frequency effects due to the leads and simplify automated assembly, although manual handling is made difficult due to their small size.
Mechanically controlled variable capacitors allow the plate spacing to be adjusted, for example by rotating or sliding a set of movable plates into alignment with a set of stationary plates. Low cost variable capacitors squeeze together alternating layers of aluminum and plastic with a screw. Electrical control of capacitance is achievable with varactors (or varicaps), which are reverse-biased semiconductor diodes whose depletion region width varies with applied voltage. They are used in phase-locked loops, amongst other applications.

Capacitor markings

Most capacitors have numbers printed on their bodies to indicate their electrical characteristics. Larger capacitors like electrolytics usually display the actual capacitance together with the unit (for example, 220 μF). Smaller capacitors like ceramics, however, use a shorthand consisting of three numbers and a letter, where the numbers show the capacitance in pF (calculated as XY x 10^Z for the numbers XYZ) and the letter indicates the tolerance (J, K or M for ±5%, ±10% and ±20% respectively).
Additionally, the capacitor may show its working voltage, temperature and other relevant characteristics.

Example

A capacitor with the text 473K 330V on its body has a capacitance of 47 x 10³ pF = 47 nF (±10%) with a working voltage of 330 V.

Applications

Main article: Applications of capacitors

Capacitors have many uses in electronic and electrical systems. They are so common that it is a rare electrical product that does not include at least one for some purpose.

Energy storage

A capacitor can store electric energy when disconnected from its charging circuit, so it can be used like a temporary battery. Capacitors are commonly used in electronic devices to maintain power supply while batteries are being changed. (This prevents loss of information in volatile memory.)
Conventional electrostatic capacitors provide less than 360 joules per kilogram of energy density, while capacitors using developing technologies can provide more than 2.52 kilo joules per kilogram^[20].
In car audio systems, large capacitors store energy for the amplifier to use on demand. Also for a flash tube a capacitor is used to hold the high voltage. In ceiling fans, capacitors play the important role of storing electrical energy to give the fan enough torque to start spinning.

Pulsed power and weapons

Groups of large, specially constructed, low-inductance high-voltage capacitors (capacitor banks) are used to supply huge pulses of current for many pulsed power applications. These include electromagnetic forming, Marx generators, pulsed lasers (especially TEA lasers), pulse forming networks, radar, fusion research, and particle accelerators.
Large capacitor banks (reservoir) are used as energy sources for the exploding-bridgewire detonators or slapper detonators in nuclear weapons and other specialty weapons. Experimental work is under way using banks of capacitors as power sources for electromagnetic armour and electromagnetic railguns and coilguns.

Power conditioning

A 10,000 microfarad capacitor in a TRM-800 amplifier

Reservoir capacitors are used in power supplies where they smooth the output of a full or half wave rectifier. They can also be used in charge pump circuits as the energy storage element in the generation of higher voltages than the input voltage.
Capacitors are connected in parallel with the power circuits of most electronic devices and larger systems (such as factories) to shunt away and conceal current fluctuations from the primary power source to provide a "clean" power supply for signal or control circuits. Audio equipment, for example, uses several capacitors in this way, to shunt away power line hum before it gets into the signal circuitry. The capacitors act as a local reserve for the DC power source, and bypass AC currents from the power supply. This is used in car audio applications, when a stiffening capacitor compensates for the inductance and resistance of the leads to the lead-acid car battery.

Power factor correction

In electric power distribution, capacitors are used for power factor correction. Such capacitors often come as three capacitors connected as a three phase load. Usually, the values of these capacitors are given not in farads but rather as a reactive power in volt-amperes reactive (VAr). The purpose is to counteract inductive loading from devices like electric motors and transmission lines to make the load appear to be mostly resistive. Individual motor or lamp loads may have capacitors for power factor correction, or larger sets of capacitors (usually with automatic switching devices) may be installed at a load center within a building or in a large utility substation.

Supression and coupling

Signal coupling

Main article: capacitive coupling

Because capacitors pass AC but block DC signals (when charged up to the applied dc voltage), they are often used to separate the AC and DC components of a signal. This method is known as AC coupling or "capacitive coupling". Here, a large value of capacitance, whose value need not be accurately controlled, but whose reactance is small at the signal frequency, is employed.

Decoupling

Main article: decoupling capacitor

A decoupling capacitor is a capacitor used to protect one part of a circuit from the effect of another, for instance to suppress noise or transients. Noise caused by other circuit elements is shunted through the capacitor, reducing the effect they have on the rest of the circuit. It is most commonly used between the power supply and ground. An alternative name is bypass capacitor as it is used to bypass the power supply or other high impedance component of a circuit.

Noise filters and snubbers

When an inductive circuit is opened, the current through the inductance collapses quickly, creating a large voltage across the open circuit of the switch or relay. If the inductance is large enough, the energy will generate a spark, causing the contact points to oxidize, deteriorate, or sometimes weld together, or destroying a solid-state switch. A snubber capacitor across the newly opened circuit creates a path for this impulse to bypass the contact points, thereby preserving their life; these were commonly found in contact breaker ignition systems, for instance. Similarly, in smaller scale circuits, the spark may not be enough to damage the switch but will still radiate undesirable radio frequency interference (RFI), which a filter capacitor absorbs. Snubber capacitors are usually employed with a low-value resistor in series, to dissipate energy and minimize RFI. Such resistor-capacitor combinations are available in a single package.
Capacitors are also used in parallel to interrupt units of a high-voltage circuit breaker in order to equally distribute the voltage between these units. In this case they are called grading capacitors.
In schematic diagrams, a capacitor used primarily for DC charge storage is often drawn vertically in circuit diagrams with the lower, more negative, plate drawn as an arc. The straight plate indicates the positive terminal of the device, if it is polarized (see electrolytic capacitor).

Motor starters

Main article: motor capacitor

In single phase squirrel cage motors, the primary winding within the motor housing is not capable of starting a rotational motion on the rotor, but is capable of sustaining one. To start the motor, a secondary winding is used in series with a non-polarized starting capacitor to introduce a lag in the sinusoidal current through the starting winding. When the secondary winding is placed at an angle with respect to the primary winding, a rotating electric field is created. The force of the rotational field is not constant, but is sufficient to start the rotor spinning. When the rotor comes close to operating speed, a centrifugal switch (or current-sensitive relay in series with the main winding) disconnects the capacitor. The start capacitor is typically mounted to the side of the motor housing. These are called capacitor-start motors, that have relatively high starting torque.
There are also capacitor-run induction motors which have a permanently-connected phase-shifting capacitor in series with a second winding. The motor is much like a two-phase induction motor.
Motor-starting capacitors are typically non-polarized electrolytic types, while running capacitors are conventional paper or plastic film dielectric types.

Signal processing

The energy stored in a capacitor can be used to represent information, either in binary form, as in DRAMs, or in analogue form, as in analog sampled filters and CCDs. Capacitors can be used in analog circuits as components of integrators or more complex filters and in negative feedback loop stabilization. Signal processing circuits also use capacitors to integrate a current signal.

Tuned circuits

Capacitors and inductors are applied together in tuned circuits to select information in particular frequency bands. For example, radio receivers rely on variable capacitors to tune the station frequency. Speakers use passive analog crossovers, and analog equalizers use capacitors to select different audio bands.
The resonant frequency f of a tuned circuit is a function of the inductance (L) and capacitance (C) in series, and is given by:

where L is in henries and C is in farads.

Sensing

Most capacitors are designed to maintain a fixed physical structure. However, various factors can change the structure of the capacitor, and the resulting change in capacitance can be used to sense those factors.
Changing the dielectric:

The effects of varying the physical and/or electrical characteristics of the dielectric can be used for sensing purposes. Capacitors with an exposed and porous dielectric can be used to measure humidity in air. Capacitors are used to accurately measure the fuel level in airplanes; as the fuel covers more of a pair of plates, the circuit capacitance increases.

Changing the distance between the plates:

Capacitors with a flexible plate can be used to measure strain or pressure. Industrial pressure transmitters used for process control use pressure-sensing diaphragms, which form a capacitor plate of an oscillator circuit. Capacitors are used as the sensor in condenser microphones, where one plate is moved by air pressure, relative to the fixed position of the other plate. Some accelerometers use MEMS capacitors etched on a chip to measure the magnitude and direction of the acceleration vector. They are used to detect changes in acceleration, e.g. as tilt sensors or to detect free fall, as sensors triggering airbag deployment, and in many other applications. Some fingerprint sensors use capacitors. Additionally, a user can adjust the pitch of a theremin musical instrument by moving his hand since this changes the effective capacitance between the user's hand and the antenna.

Changing the effective area of the plates:

Capacitive touch switches are now used on many consumer electronic products.

Hazards and safety

Capacitors may retain a charge long after power is removed from a circuit; this charge can cause dangerous or even potentially fatal shocks or damage connected equipment. For example, even a seemingly innocuous device such as a disposable camera flash unit powered by a 1.5 volt AA battery contains a capacitor which may be charged to over 300 volts. This is easily capable of delivering a shock. Service procedures for electronic devices usually include instructions to discharge large or high-voltage capacitors. Capacitors may also have built-in discharge resistors to dissipate stored energy to a safe level within a few seconds after power is removed. High-voltage capacitors are stored with the terminals shorted, as protection from potentially dangerous voltages due to dielectric absorption.
Some old, large oil-filled capacitors contain polychlorinated biphenyls (PCBs). It is known that waste PCBs can leak into groundwater under landfills. Capacitors containing PCB were labelled as containing "Askarel" and several other trade names. PCB-filled capacitors are found in very old (pre 1975) fluorescent lamp ballasts, and other applications.
High-voltage capacitors may catastrophically fail when subjected to voltages or currents beyond their rating, or as they reach their normal end of life. Dielectric or metal interconnection failures may create arcing that vaporizes dielectric fluid, resulting in case bulging, rupture, or even an explosion. Capacitors used in RF or sustained high-current applications can overheat, especially in the center of the capacitor rolls. Capacitors used within high-energy capacitor banks can violently explode when a short in one capacitor causes sudden dumping of energy stored in the rest of the bank into the failing unit. High voltage vacuum capacitors can generate soft X-rays even during normal operation. Proper containment, fusing, and preventive maintenance can help to minimize these hazards.
High-voltage capacitors can benefit from a pre-charge to limit in-rush currents at power-up of high voltage direct current (HVDC) circuits. This will extend the life of the component and may mitigate high-voltage hazards.

Resistor

2010-05-16T20:58:00.001-07:00

A resistor is a two-terminal electronic component that produces a voltage across its terminals that is proportional to the electric current passing through it in accordance with Ohm's law:

V = I R

Resistors are elements of electrical networks and electronic circuits and are ubiquitous in most electronic equipment. Practical resistors can be made of various compounds and films, as well as resistance wire (wire made of a high-resistivity alloy, such as nickel/chrome).
The primary characteristics of a resistor are the resistance, the tolerance, maximum working voltage and the power rating. Other characteristics include temperature coefficient, noise, and inductance. Less well-known is critical resistance, the value below which power dissipation limits the maximum permitted current flow, and above which the limit is applied voltage. Critical resistance is determined by the design, materials and dimensions of the resistor.
Resistors can be integrated into hybrid and printed circuits, as well as integrated circuits. Size, and position of leads (or terminals) are relevant to equipment designers; resistors must be physically large enough not to overheat when dissipating their power.

Units

The ohm (symbol: Ω) is the SI unit of electrical resistance, named after Georg Simon Ohm. Commonly used multiples and submultiples in electrical and electronic usage are the milliohm (1x10⁻³), kilohm (1x10³), and megohm (1x10⁶).

Theory of operation

Ohm's law

The behavior of an ideal resistor is dictated by the relationship specified in Ohm's law:

Ohm's law states that the voltage (V) across a resistor is proportional to the current (I) through it where the constant of proportionality is the resistance (R).
Equivalently, Ohm's law can be stated:

This formulation of Ohm's law states that, when a voltage (V) is maintained across a resistance (R), a current (I) will flow through the resistance.
This formulation is often used in practice. For example, if V is 12 volts and R is 400 ohms, a current of 12 / 400 = 0.03 amperes will flow through the resistance R.

Series and parallel resistors

Main article: Series and parallel circuits

Resistors in a parallel configuration each have the same potential difference (voltage). To find their total equivalent resistance (R_eq):

The parallel property can be represented in equations by two vertical lines "||" (as in geometry) to simplify equations. For two resistors,

The current through resistors in series stays the same, but the voltage across each resistor can be different. The sum of the potential differences (voltage) is equal to the total voltage. To find their total resistance:

A resistor network that is a combination of parallel and series can be broken up into smaller parts that are either one or the other. For instance,

However, many resistor networks cannot be split up in this way. Consider a cube, each edge of which has been replaced by a resistor. For example, determining the resistance between two opposite vertices requires additional transforms, such as the Y-Δ transform, or else matrix methods must be used for the general case. However, if all twelve resistors are equal, the corner-to-corner resistance is ⁵⁄₆ of any one of them.
The practical application to resistors is that a resistance of any non-standard value can be obtained by connecting standard values in series or in parallel.

Power dissipation

The power dissipated by a resistor (or the equivalent resistance of a resistor network) is calculated using the following:
All three equations are equivalent. The first is derived from Joule's first law. Ohm’s Law derives the other two from that.
The total amount of heat energy released is the integral of the power over time:

If the average power dissipated is more than the resistor can safely dissipate, the resistor may depart from its nominal resistance and may become damaged by overheating. Excessive power dissipation may raise the temperature of the resistor to a point where it burns out, which could cause a fire in adjacent components and materials. There are flameproof resistors that fail (open circuit) before they overheat dangerously.
Note that the nominal power rating of a resistor is not the same as the power that it can safely dissipate in practical use. Air circulation and proximity to a circuit board, ambient temperature, and other factors can reduce acceptable dissipation significantly. Rated power dissipation may be given for an ambient temperature of 25 °C in free air. Inside an equipment case at 60 °C, rated dissipation will be significantly less; if we are dissipating a bit less than the maximum figure given by the manufacturer we may still be outside the safe operating area, and courting premature failure.

Construction

A single in line (SIL) resistor package with 8 individual, 47 ohm resistors. One end of each resistor is connected to a separate pin and the other ends are all connected together to the remaining (common) pin - pin 1, at the end identified by the white dot.

Lead arrangements

Through-hole components typically have leads leaving the body axially. Others have leads coming off their body radially instead of parallel to the resistor axis. Other components may be SMT (surface mount technology) while high power resistors may have one of their leads designed into the heat sink.

Carbon composition

Carbon composition resistors consist of a solid cylindrical resistive element with embedded wire leads or metal end caps to which the lead wires are attached. The body of the resistor is protected with paint or plastic. Early 20th-century carbon composition resistors had uninsulated bodies; the lead wires were wrapped around the ends of the resistance element rod and soldered. The completed resistor was painted for color coding of its value.
The resistive element is made from a mixture of finely ground (powdered) carbon and an insulating material (usually ceramic). A resin holds the mixture together. The resistance is determined by the ratio of the fill material (the powdered ceramic) to the carbon. Higher concentrations of carbon, a weak conductor, result in lower resistance. Carbon composition resistors were commonly used in the 1960s and earlier, but are not so popular for general use now as other types have better specifications, such as tolerance, voltage dependence, and stress (carbon composition resistors will change value when stressed with over-voltages). Moreover, if internal moisture content (from exposure for some length of time to a humid environment) is significant, soldering heat will create a non-reversible change in resistance value. These resistors, however, if never subjected to overvoltage nor overheating were remarkably reliable considering the components size ^[1]
They are still available, but comparatively quite costly. Values ranged from fractions of an ohm to 22 megohms. Because of the expense these resistors tend to be used in power supplies and welding controls^[1].

Carbon film

A carbon film is deposited on an insulating substrate, and a helix cut in it to create a long, narrow resistive path. Varying shapes, coupled with the resistivity of carbon, (ranging from 90 to 400 nΩm) can provide a variety of resistances.^[2] Carbon film resistors feature a power rating range of 0.125 W to 5 W at 70 °C. Resistances available range from 1 ohm to 10 megohm. The carbon film resistor has an operating temperature range of -55 °C to 155 °C. It has 200 to 600 volts maximum working voltage range.^[3]. Special carbon film resistors are used in applications requiring high pulse stability^[1].

Thick and thin film

Thick film resistors became popular during the 1970s, and most SMD (surface mount device) resistors today are of this type. The principal difference between thin film and thick film resistors is not the actual thickness of the film, but rather how the film is applied to the cylinder (axial resistors) or the surface (SMD resistors).
Thin film resistors are made by sputtering (a method of vacuum deposition) the resistive material onto an insulating substrate. The film is then etched in a similar manner to the old (subtractive) process for making printed circuit boards; that is, the surface is coated with a photo-sensitive material, then covered by a pattern film, irradiated with ultraviolet light, and then the exposed photo-sensitive coating is developed, and underlying thin film is etched away.
Thick film resistors are manufactured using screen and stencil printing processes^[1].
Because the time during which the sputtering is performed can be controlled, the thickness of the thin film can be accurately controlled. The type of material is also usually different consisting of one or more ceramic (cermet) conductors such as tantalum nitride (TaN), ruthenium dioxide (RuO₂), lead oxide (PbO), bismuth ruthenate (Bi₂Ru₂O₇), nickel chromium (NiCr), and/or bismuth iridate (Bi₂Ir₂O₇).
The resistance of both thin and thick film resistors after manufacture is not highly accurate; they are usually trimmed to an accurate value by abrasive or laser trimming. Thin film resistors are usually specified with tolerances of 0.1, 0.2, 0.5, or 1%, and with temperature coefficients of 5 to 25 ppm/K.
Thick film resistors may use the same conductive ceramics, but they are mixed with sintered (powdered) glass and some kind of liquid so that the composite can be screen-printed. This composite of glass and conductive ceramic (cermet) material is then fused (baked) in an oven at about 850 °C.
Thick film resistors, when first manufactured, had tolerances of 5%, but standard tolerances have improved to 2% or 1% in the last few decades. Temperature coefficients of thick film resistors are high, typically ±200 or ±250 ppm/K; a 40 kelvin (70 °F) temperature change can change the resistance by 1%.
Thin film resistors are usually far more expensive than thick film resistors. For example, SMD thin film resistors, with 0.5% tolerances, and with 25 ppm/K temperature coefficients, when bought in full size reel quantities, are about twice the cost of 1%, 250 ppm/K thick film resistors.

Metal film

A common type of axial resistor today is referred to as a metal-film resistor. Metal electrode leadless face (MELF) resistors often use the same technology, but are a cylindrically shaped resistor designed for surface mounting. Note that other types of resistors (e.g., carbon composition) are also available in MELF packages.
Metal film resistors are usually coated with nickel chromium (NiCr), but might be coated with any of the cermet materials listed above for thin film resistors. Unlike thin film resistors, the material may be applied using different techniques than sputtering (though that is one such technique). Also, unlike thin-film resistors, the resistance value is determined by cutting a helix through the coating rather than by etching. (This is similar to the way carbon resistors are made.) The result is a reasonable tolerance (0.5, 1, or 2%) and a temperature coefficient that is generally between 50 and 100 ppm/K.^[4]. Metal film resistors process good noise characteristics and low non-linearity. Also beneficial are the components efficient tolerance, temperature coefficient and stability^[1].

Wirewound

Types of windings in wire resistors:
1 - common
2 - bifilar
3 - common on a thin former
4 - Ayrton-Perry

Wirewound resistors are commonly made by winding a metal wire, usually nichrome, around a ceramic, plastic, or fiberglass core. The ends of the wire are soldered or welded to two caps or rings, attached to the ends of the core. The assembly is protected with a layer of paint, molded plastic, or an enamel coating baked at high temperature. Because of the very high surface temperature these resistors can withstand temperatures of up to +450°C^[1]. Wire leads in low power wirewound resistors are usually between 0.6 and 0.8 mm in diameter and tinned for ease of soldering. For higher power wirewound resistors, either a ceramic outer case or an aluminum outer case on top of an insulating layer is used. The aluminum-cased types are designed to be attached to a heat sink to dissipate the heat; the rated power is dependent on being used with a suitable heat sink, e.g., a 50 W power rated resistor will overheat at a fraction of the power dissipation if not used with a heat sink. Large wirewound resistors may be rated for 1,000 watts or more.
Because wirewound resistors are coils they have more undesirable inductance than other types of resistor, although winding the wire in sections with alternately reversed direction can minimize inductance. Other techniques employ bifilar winding, or a flat thin former (to reduce cross-section area of the coil). For most demanding circuits resistors with Ayrton-Perry winding are used.
Applications of wirewound resistors are similar to those of composition resistors with the exception of the high frequency. The high frequency of wirewound resistors is substantially worse than that of a composition resistor^[1]

Foil resistor

The primary resistance element of a foil resistor is a special alloy foil several micrometres thick. Since their introduction in the 1960s, foil resistors have had the best precision and stability of any resistor available. One of the important parameters influencing stability is the temperature coefficient of resistance (TCR). The TCR of foil resistors is extremely low, and has been further improved over the years. One range of ultra-precision foil resistors offers a TCR of 0.14 ppm/°C, tolerance ±0.005%, long-term stability (1 year) 25 ppm, (3 year) 50 ppm (further improved 5-fold by hermetic sealing), stability under load (2000 hours) 0.03%, thermal EMF 0.1 μV/°C, noise -42 dB, voltage coefficient 0.1 ppm/V, inductance 0.08 μH, capacitance 0.5 pF.^[5]

Ammeter shunts

An ammeter shunt is a special type of current-sensing resistor, having four terminals and a value in milliohms or even micro-ohms. Current-measuring instruments, by themselves, can usually accept only limited currents. To measure high currents, the current passes through the shunt, where the voltage drop is measured and interpreted as current. A typical shunt consists of two solid metal blocks, sometimes brass, mounted on to an insulating base. Between the blocks, and soldered or brazed to them, are one or more strips of low temperature coefficient of resistance (TCR) manganin alloy. Large bolts threaded into the blocks make the current connections, while much-smaller screws provide voltage connections. Shunts are rated by full-scale current, and often have a voltage drop of 50 mV at rated current. Such meters are adapted to the shunt full current rating by using an appropriately marked dial face; no change need be made to the other parts of the meter.

Grid resistor

In heavy-duty industrial high-current applications, a grid resistor is a large convection-cooled lattice of stamped metal alloy strips connected in rows between two electrodes. Such industrial grade resistors can be as large as a refrigerator; some designs can handle over 500 amperes of current, with a range of resistances extending lower than 0.04 ohms. They are used in applications such as dynamic braking and load banking for locomotives and trams, neutral grounding for industrial AC distribution, control loads for cranes and heavy equipment, load testing of generators and harmonic filtering for electric substations.^[6]^[7]^[8]
The term grid resistor is sometimes used to describe a resistor of any type connected to the control grid of a vacuum tube. This is not a resistor technology; it is an electronic circuit topology.

Negative resistors

Main article: Negative resistance

A device exhibiting negative resistance over part of its characteristic curve can be made with active circuit components.

Special varieties

Adjustable resistors

Tapped resistors

A resistor may have one or more fixed tapping points so that the resistance can be changed by moving the connecting wires to different terminals. Wire-wound power resistors can have a tapping point that can slide the resistance element, allowing any part of the resistance to be used.
Where continuous adjustment of the resistance value during operation of equipment is required, the sliding resistance tap can be connected to a knob accessible to an operator. Such a device is called a "rheostat" and has two terminals.
A frequent element of electronic devices is a three-terminal resistor with a continuously-adjustable tapping point controlled by rotation of a shaft or knob. These variable resistors are known as a potentiometer when all three terminals are connected, since they act as a continuously adjustable voltage divider. A common example is a volume control for a radio receiver.^[9]
Accurate, high-resolution panel-mounted pots have resistance elements typically wirewound on a helical mandrel, although some include a conductive-plastic resistance coating over the wire to improve resolution. These typically offer ten turns of their shafts to cover their full range. They are usually set with dials that include a simple turns counter and a graduated dial. Electronic analog computers used them in quantity for setting coefficients, and delayed-sweep oscilloscopes of recent decades included one on their panels.

Strain gauges

Main article: Strain gauge

The strain gauge, invented by Edward E. Simmons and Arthur C. Ruge in 1938, is a type of resistor that changes value with applied strain. A single resistor may be used, or a pair (half bridge), or four resistors connected in a Wheatstone bridge configuration. The strain resistor is bonded with adhesive to an object that will be subjected to mechanical strain. With the strain gauge and a filter, amplifier, and analog/digital converter, the strain on an object can be measured.

Resistance decade boxes

A resistance decade box is a box containing resistors of many values and two (or four) terminals, with a mechanical switch that allows a resistance of any value allowed by the box to be dialed. Usually the resistance is accurate to high precision, ranging from laboratory/calibration grade accurate to within 20 parts per million, to field grade at 1%. Inexpensive boxes with lesser accuracy are also available. All types offer a convenient way of selecting and quickly changing a resistance in laboratory, experimental and development work without having to stock and seek individual resistors of the required value. The range of resistance provided, the maximum resolution, and the accuracy characterize the box. For example, one box offers resistances from 0 to 24 megohms, maximum resolution 0.1 ohm, accuracy 0.1%.^[10]

Special varieties

There are special types of resistor whose resistance changes with various quantities: the resistance of thermistors varies greatly with temperature, whether external or due to dissipation, so they can be used for temperature or current sensing; metal oxide varistors drop to a very low resistance when a high voltage is applied, making them suitable for over-voltage protection; the resistance of photoresistors varies with illumination; the resistance of a Quantum Tunnelling Composite can vary by a factor of 10¹² with mechanical pressure applied; and so on.

Measurement

The value of a resistor can be measured with an ohmmeter, which may be one function of a multimeter. Usually, probes on the ends of test leads connect to the resistor.
Measuring low-value resistors, such as fractional-ohm resistors, with acceptable accuracy requires four-terminal connections. One pair of terminals applies a known, calibrated current to the resistor, while the other pair senses the voltage drop across the resistor. Some laboratory test instruments have spring-loaded pairs of contacts, with neighboring contacts electrically isolated from each other. Better digital multimeters have four terminals on their panels, generally used with special test leads. These comprise four wires in all, and have special test clips with jaws insulated from each other. One jaw provides the measuring current, while the other senses the voltage drop. The resistance is then calculated using Ohm's Law.

Standards

Production resistors

Resistor characteristics are quantified and reported using various national standards. In the US, MIL-STD-202^[11] contains the relevant test methods to which other standards refer.
There are various standards specifying properties of resistors for use in equipment:

BS 1852
EIA-RS-279
MIL-PRF-26
MIL-PRF-39007 (Fixed Power, established reliability)
MIL-PRF-55342 (Surface-mount thick and thin film)
MIL-PRF-914
MIL-R-11
MIL-R-39017 (Fixed, General Purpose, Established Reliability)
MIL-PRF-32159 (zero ohm jumpers)

There are other United States military procurement MIL-R- standards.

Resistance standards

The primary standard for resistance, the "mercury ohm" was initially defined in 1884 in as a column of mercury 106mm long and 1 square millimeter in cross-section, at 0 degrees Celsius. Difficulties in precisely measuring the physical constants to replicate this standard result in variations of as much as 30 ppm. From 1900 the mercury ohm was replaced with a precision machined plate of manganin^[12]. Since 1990 the international resistance standard has been based on the quantized Hall effect discovered by Klaus von Klitzing, for which he won the Nobel Prize in Physics in 1985^[13]
Resistors of extremely high precision are manufactured for calibration and laboratory use. They may have four terminals, using one pair to carry an operating current and the other pair to measure the voltage drop; this eliminates errors caused by voltage drops across the lead resistances, because no current flows through voltage sensing leads. It is important in small value resistors (100–0.0001 Ohm) where lead resistance is significant or even comparable with respect to resistance standard value.^[14]

Resistor marking

Main article: Electronic color code

Most axial resistors use a pattern of colored stripes to indicate resistance. Surface-mount resistors are marked numerically, if they are big enough to permit marking; more-recent small sizes are impractical to mark. Cases are usually tan, brown, blue, or green, though other colors are occasionally found such as dark red or dark gray.
Early 20th century resistors, essentially uninsulated, were dipped in paint to cover their entire body for color coding. A second color of paint was applied to one end of the element, and a color dot (or band) in the middle provided the third digit. The rule was "body, tip, dot", providing two significant digits for value and the decimal multiplier, in that sequence. Default tolerance was ±20%. Closer-tolerance resistors had silver (±10%) or gold-colored (±5%) paint on the other end.

Four-band resistors

Four-band identification is the most commonly used color-coding scheme on resistors. It consists of four colored bands that are painted around the body of the resistor. The first two bands encode the first two significant digits of the resistance value, the third is a power-of-ten multiplier or number-of-zeroes, and the fourth is the tolerance accuracy, or acceptable error, of the value. The first three bands are equally spaced along the resistor; the spacing to the fourth band is wider. Sometimes a fifth band identifies the thermal coefficient, but this must be distinguished from the true 5-color system, with 3 significant digits.
For example, green-blue-yellow-red is 56×10⁴ Ω = 560 kΩ ± 2%. An easier description can be as followed: the first band, green, has a value of 5 and the second band, blue, has a value of 6, and is counted as 56. The third band, yellow, has a value of 10⁴, which adds four 0's to the end, creating 560,000 Ω at ±2% tolerance accuracy. 560,000 Ω changes to 560 kΩ ±2% (as a kilo- is 10³).
Each color corresponds to a certain digit, progressing from darker to lighter colors, as shown in the chart below.

Color	1^st band	2^nd band	3^rd band (multiplier)	4^th band (tolerance)	Temp. Coefficient
Black	0	0	×10⁰
Brown	1	1	×10¹	±1% (F)	100 ppm
Red	2	2	×10²	±2% (G)	50 ppm
Orange	3	3	×10³		15 ppm
Yellow	4	4	×10⁴		25 ppm
Green	5	5	×10⁵	±0.5% (D)
Blue	6	6	×10⁶	±0.25% (C)
Violet	7	7	×10⁷	±0.1% (B)
Gray	8	8	×10⁸	±0.05% (A)
White	9	9	×10⁹
Gold			×10⁻¹	±5% (J)
Silver			×10⁻²	±10% (K)
None				±20% (M)

There are many mnemonics for remembering these colors.

Preferred values

Main article: Preferred number

Early resistors were made in more or less arbitrary round numbers; a series might have 100, 125, 150, 200, 300, etc. Resistors as manufactured are subject to a certain percentage tolerance, and it makes sense to manufacture values that correlate with the tolerance, so that the actual value of a resistor overlaps slightly with its neighbors. Wider spacing leaves gaps; narrower spacing increases manufacturing and inventory costs to provide resistors that are more or less interchangeable.
A logical scheme is to produce resistors in a range of values which increase in a geometrical progression, so that each value is greater than its predecessor by a fixed multiplier or percentage, chosen to match the tolerance of the range. For example, for a tolerance of ±20% it makes sense to have each resistor about 1.5 times its predecessor, covering a decade in 6 values. In practice the factor used is 1.4678, giving values of 1.47, 2.15, 3.16, 4.64, 6.81, 10 for the 1-10 decade (a decade is a range increasing by a factor of 10; 0.1-1 and 10-100 are other examples); these are rounded in practice to 1.5, 2.2, 3.3, 4.7, 6.8, 10; followed, of course by 15, 22, 33, … and preceded by … 0.47, 0.68, 1. This scheme has been adopted as the E6 range of the IEC 60063 preferred number series. There are also E12, E24, E48, E96 and E192 ranges for components of ever tighter tolerance, with 12, 24, 96, and 192 different values within each decade. The actual values used are in the IEC 60063 lists of preferred numbers.
A resistor of 100 ohms ±20% would be expected to have a value between 80 and 120 ohms; its E6 neighbors are 68 (54-82) and 150 (120-180) ohms. A sensible spacing, E6 is used for ±20% components; E12 for ±10%; E24 for ±5%; E48 for ±2%, E96 for ±1%; E192 for ±0.5% or better. Resistors are manufactured in values from a few milliohms to about a gigaohm in IEC60063 ranges appropriate for their tolerance.
Earlier power wirewound resistors, such as brown vitreous-enameled types, however, were made with a different system of preferred values, such as some of those mentioned in the first sentence of this section.

5-band axial resistors

5-band identification is used for higher precision (lower tolerance) resistors (1%, 0.5%, 0.25%, 0.1%), to specify a third significant digit. The first three bands represent the significant digits, the fourth is the multiplier, and the fifth is the tolerance. Five-band resistors with a gold or silver 4th band are sometimes encountered, generally on older or specialized resistors. The 4th band is the tolerance and the 5th the temperature coefficient.

SMD resistors

This image shows four surface-mount resistors (the component at the upper left is a capacitor) including two zero-ohm resistors. Zero-ohm links are often used instead of wire links, so that they can be inserted by a resistor-inserting machine. Of course, their resistance is finite, although quite low. Zero is simply a brief description of their function.

Surface mounted resistors are printed with numerical values in a code related to that used on axial resistors. Standard-tolerance surface-mount technology (SMT) resistors are marked with a three-digit code, in which the first two digits are the first two significant digits of the value and the third digit is the power of ten (the number of zeroes). For example:

334	= 33 × 10,000 ohms = 330 kilohms
222	= 22 × 100 ohms = 2.2 kilohms
473	= 47 × 1,000 ohms = 47 kilohms
105	= 10 × 100,000 ohms = 1.0 megohm

Resistances less than 100 ohms are written: 100, 220, 470. The final zero represents ten to the power zero, which is 1. For example:

100	= 10 × 1 ohm = 10 ohms
220	= 22 × 1 ohm = 22 ohms

Sometimes these values are marked as 10 or 22 to prevent a mistake.
Resistances less than 10 ohms have 'R' to indicate the position of the decimal point (radix point). For example:

4R7	= 4.7 ohms
0R22	= 0.22 ohms
0R01	= 0.01 ohms

Precision resistors are marked with a four-digit code, in which the first three digits are the significant figures and the fourth is the power of ten. For example:

1001	= 100 × 10 ohms = 1.00 kilohm
4992	= 499 × 100 ohms = 49.9 kilohm
1000	= 100 × 1 ohm = 100 ohms

000 and 0000 sometimes appear as values on surface-mount zero-ohm links, since these have (approximately) zero resistance.
More recent surface-mount resistors are too small, physically, to permit practical markings to be applied.

Industrial type designation

Format: [two letters][resistance value (three digit)][tolerance code(numerical - one digit)] ^[15]

Power Rating at 70 °C
Type No.	Power rating (watts)	MIL-R-11 Style	MIL-R-39008 Style
BB	⅛	RC05	RCR05
CB	¼	RC07	RCR07
EB	½	RC20	RCR20
GB	1	RC32	RCR32
HB	2	RC42	RCR42
GM	3	-	-
HM	4	-	-

Tolerance Code
Industrial type designation	Tolerance	MIL Designation
5	±5%	J
2	±20%	M
1	±10%	K
-	±2%	G
-	±1%	F
-	±0.5%	D
-	±0.25%	C
-	±0.1%	B

The operational temperature range distinguishes commercial grade, industrial grade and military grade components.

Commercial grade: 0 °C to 70 °C
Industrial grade: −40 °C to 85 °C (sometimes −25 °C to 85 °C)
Military grade: −55 °C to 125 °C (sometimes -65 °C to 275 °C)
Standard Grade -5 °C to 60 °C

Electrical and thermal noise

In precision applications it is often necessary to minimize electronic noise. As dissipative elements, even ideal resistors will naturally produce a fluctuating "noise" voltage across their terminals. This Johnson–Nyquist noise is a fundamental noise source which depends only upon the temperature and resistance of the resistor, and is predicted by the fluctuation–dissipation theorem. For example, the gain in a simple (non-) inverting amplifier is set using a voltage divider. Noise considerations dictate that the smallest practical resistance should be used, since the Johnson–Nyquist noise voltage scales with resistance, and any resistor noise in the voltage divider will be impressed upon the amplifier's output.
In addition, small voltage differentials may appear on the resistors due to thermoelectric effect if their ends are not kept at the same temperature. The voltages appear in the junctions of the resistor leads with the circuit board and with the resistor body. Common metal film resistors show such an effect at a magnitude of about 20 µV/°C. Some carbon composition resistors can go as high as 400 µV/°C, and specially constructed resistors can go as low as 0.05 µV/°C. In applications where thermoelectric effects may become important, care has to be taken (for example) to mount the resistors horizontally to avoid temperature gradients and to mind the air flow over the board.^[16]
Practical resistors frequently exhibit other, "non-fundamental", sources of noise, usually called "excess noise." Excess noise results in a "Noise Index" for a type of resistor. Excess Noise is due to current flow in the resistor and is specified as μV/V/decade - μV of noise per volt applied across the resistor per decade of frequency. The μV/V/decade value is frequently given in dB so that a resistor with a noise index of 0dB will exhibit 1 μV (rms) of excess noise for each volt across the resistor in each frequency decade. Excess noise is an example of 1/f noise. Thick-film and carbon composition resistors generate more noise than other types at low frequencies; wire-wound and thin-film resistors, though much more expensive, are often utilized for their better noise characteristics. Carbon composition resistors can exhibit a noise index of 0 dB while bulk metal foil resistors may have a noise index of -40 dB, usually making the excess noise of metal foil resistors insignificant.^[17]
Thin film surface mount resistors typically have lower noise and better thermal stability than thick film surface mount resistors. However, the design engineer must read the data sheets for the family of devices to weigh the various device tradeoffs.

Failure modes

Like every part, resistors can fail in normal use. Thermal and mechanical stress, humidity, etc., can play a part. Carbon composition resistors and metal film resistors typically fail as open circuits. Carbon-film resistors may decrease or increase in resistance.^[18] Carbon film and composition resistors can open if running close to their maximum dissipation. This is also possible but less likely with metal film and wirewound resistors. If not enclosed, wirewound resistors can corrode. The resistance of carbon composition resistors are prone to drift over time and are easily damaged by excessive heat in soldering (the binder evaporates). Variable resistors become electrically noisy as they wear.
All resistors can be destroyed, usually by going open-circuit, if subjected to excessive current due to failure of other components or accident.

Diode

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Diode

In electronics, a diode is a two-terminal electronic component that conducts electric current in only one direction. The term usually refers to a semiconductor diode, the most common type today, which is a crystal of semiconductor connected to two electrical terminals. A vacuum tube diode (now little used) is a vacuum tube with two electrodes; a plate and a cathode.
The most common function of a diode is to allow an electric current to flow though it in one direction (called the diode's forward direction) while blocking current in the opposite direction (the reverse direction). Thus, the diode can be thought of as an electronic version of a check valve. This unidirectional behavior is called rectification, and is used to convert alternating current to direct current, and extract modulation from radio signals in radio receivers.
However, diodes can have more complicated behavior than this simple on-off action, due to their complex non-linear electrical characteristics, which can be tailored by varying the construction of their P-N junction. These are exploited in special purpose diodes that perform many different functions. Diodes are used to regulate voltage (Zener diodes), electronically tune radio and TV receivers (varactor diodes), generate radio frequency oscillations (tunnel diodes), and produce light (light emitting diodes).
Diodes were the first semiconductor electronic devices. The discovery of crystals' rectifying abilities was made by German physicist Ferdinand Braun in 1874. The first semiconductor diodes, called cat's whisker diodes were made of crystals of minerals such as galena. Today most diodes are made of silicon, but other semiconductors such as germanium are sometimes used.

History

Although the crystal semiconductor diode was popular before the thermionic diode, thermionic and solid state diodes were developed in parallel.
The basic principle of operation of thermionic diodes was discovered by Frederick Guthrie in 1873.^[1] Guthrie discovered that a positively-charged electroscope could be discharged by bringing a grounded piece of white-hot metal close to it (but not actually touching it). The same did not apply to a negatively charged electroscope, indicating that the current flow was only possible in one direction.
The principle was independently rediscovered by Thomas Edison on February 13, 1880. At the time Edison was carrying out research into why the filaments of his carbon-filament light bulbs nearly always burned out at the positive-connected end. He had a special bulb made with a metal plate sealed into the glass envelope, and he was able to confirm that an invisible current could be drawn from the glowing filament through the vacuum to the metal plate, but only when the plate was connected to the positive supply.
Edison devised a circuit where his modified light bulb more or less replaced the resistor in a DC voltmeter and on this basis was awarded a patent for it in 1883.^[2] There was no apparent practical use for such device at the time, and the patent application was most likely simply a precaution in case someone else did find a use for the so-called “Edison Effect”.
About 20 years later, John Ambrose Fleming (scientific adviser to the Marconi Company and former Edison employee) realized that the Edison effect could be used as a precision radio detector. Fleming patented the first true thermionic diode in Britain ^[3] on November 16, 1904 (followed by U.S. Patent 803,684 in November 1905).
The principle of operation of crystal diodes was discovered in 1874 by the German scientist Karl Ferdinand Braun.^[4] Braun patented the crystal rectifier in 1899.^[5] Braun’s discovery was further developed by Jagdish Chandra Bose into a useful device for radio detection.
The first actual radio receiver using a crystal diode was built by Greenleaf Whittier Pickard. Pickard received a patent for a silicon crystal detector on November 20, 1906.^[6]
Other experimenters tried a variety of minerals and other substances, although by far the most popular was the lead sulfide mineral Galena. Although other substances offered slightly better performance, galena had the advantage of being cheap and easy to obtain, and was used almost exclusively in home-built “crystal sets”, until the advent of inexpensive fixed-germanium diodes in the 1950s.
At the time of their invention, such devices were known as rectifiers. In 1919, William Henry Eccles coined the term diode from the Greek roots dia, meaning “through”, and ode (from ὅδος), meaning “path”.

Thermionic and gaseous state diodes

Figure 4: The symbol for an indirect heated vacuum tube diode. From top to bottom, the components are the anode, the cathode, and the heater filament.

Thermionic diodes are thermionic-valve devices (also known as vacuum tubes, tubes, or valves), which are arrangements of electrodes surrounded by a vacuum within a glass envelope. Early examples were fairly similar in appearance to incandescent light bulbs.
In thermionic valve diodes, a current through the heater filament indirectly heats the cathode, another internal electrode treated with a mixture of barium and strontium oxides, which are oxides of alkaline earth metals; these substances are chosen because they have a small work function. (Some valves use direct heating, in which a tungsten filament acts as both heater and cathode.) The heat causes thermionic emission of electrons into the vacuum. In forward operation, a surrounding metal electrode called the anode is positively charged so that it electrostatically attracts the emitted electrons. However, electrons are not easily released from the unheated anode surface when the voltage polarity is reversed. Hence, any reverse flow is negligible.
For much of the 20th century, thermionic valve diodes were used in analog signal applications, and as rectifiers in many power supplies. Today, valve diodes are only used in niche applications such as rectifiers in electric guitar and high-end audio amplifiers as well as specialized high-voltage equipment.

Semiconductor diodes

A modern semiconductor diode is made of a crystal of semiconductor like silicon that has impurities added to it to create a region on one side that contains negative charge carriers (electrons), called n-type semiconductor, and a region on the other side that contains positive charge carriers (holes), called p-type semiconductor. The diode's terminals are attached to each of these regions. The boundary within the crystal between these two regions, called a PN junction, is where the action of the diode takes place. The crystal conducts conventional current in a direction from the p-type side (called the anode) to the n-type side (called the cathode), but not in the opposite direction.
Another type of semiconductor diode, the Schottky diode, is formed from the contact between a metal and a semiconductor rather than by a p-n junction.

Current–voltage characteristic

A semiconductor diode’s behavior in a circuit is given by its current–voltage characteristic, or I–V graph (see graph at right). The shape of the curve is determined by the transport of charge carriers through the so-called depletion layer or depletion region that exists at the p-n junction between differing semiconductors. When a p-n junction is first created, conduction band (mobile) electrons from the N-doped region diffuse into the P-doped region where there is a large population of holes (vacant places for electrons) with which the electrons “recombine”. When a mobile electron recombines with a hole, both hole and electron vanish, leaving behind an immobile positively charged donor (dopant) on the N-side and negatively charged acceptor (dopant) on the P-side. The region around the p-n junction becomes depleted of charge carriers and thus behaves as an insulator.
However, the width of the depletion region (called the depletion width) cannot grow without limit. For each electron-hole pair that recombines, a positively-charged dopant ion is left behind in the N-doped region, and a negatively charged dopant ion is left behind in the P-doped region. As recombination proceeds more ions are created, an increasing electric field develops through the depletion zone which acts to slow and then finally stop recombination. At this point, there is a “built-in” potential across the depletion zone.
If an external voltage is placed across the diode with the same polarity as the built-in potential, the depletion zone continues to act as an insulator, preventing any significant electric current flow (unless electron/hole pairs are actively being created in the junction by, for instance, light. see photodiode). This is the reverse bias phenomenon. However, if the polarity of the external voltage opposes the built-in potential, recombination can once again proceed, resulting in substantial electric current through the p-n junction (i.e. substantial numbers of electrons and holes recombine at the junction). For silicon diodes, the built-in potential is approximately 0.7 V. Thus, if an external current is passed through the diode, about 0.7 V will be developed across the diode such that the P-doped region is positive with respect to the N-doped region and the diode is said to be “turned on” as it has a forward bias.

Figure 5: I–V characteristics of a P-N junction diode (not to scale).

A diode’s 'I–V characteristic' can be approximated by four regions of operation (see the figure at right).
At very large reverse bias, beyond the peak inverse voltage or PIV, a process called reverse breakdown occurs which causes a large increase in current (i.e. a large number of electrons and holes are created at, and move away from the pn junction) that usually damages the device permanently. The avalanche diode is deliberately designed for use in the avalanche region. In the zener diode, the concept of PIV is not applicable. A zener diode contains a heavily doped p-n junction allowing electrons to tunnel from the valence band of the p-type material to the conduction band of the n-type material, such that the reverse voltage is “clamped” to a known value (called the zener voltage), and avalanche does not occur. Both devices, however, do have a limit to the maximum current and power in the clamped reverse voltage region. Also, following the end of forward conduction in any diode, there is reverse current for a short time. The device does not attain its full blocking capability until the reverse current ceases.
The second region, at reverse biases more positive than the PIV, has only a very small reverse saturation current. In the reverse bias region for a normal P-N rectifier diode, the current through the device is very low (in the µA range). However, this is temperature dependent, and at suffiently high temperatures, a substantial amount of reverse current can be observed (mA or more).
The third region is forward but small bias, where only a small forward current is conducted.
As the potential difference is increased above an arbitrarily defined “cut-in voltage” or “on-voltage” or “diode forward voltage drop (V_d)”, the diode current becomes appreciable (the level of current considered “appreciable” and the value of cut-in voltage depends on the application), and the diode presents a very low resistance.
The current–voltage curve is exponential. In a normal silicon diode at rated currents, the arbitrary “cut-in” voltage is defined as 0.6 to 0.7 volts. The value is different for other diode types — Schottky diodes can be rated as low as 0.2 V and red or blue light-emitting diodes (LEDs) can have values of 1.4 V and 4.0 V respectively.
At higher currents the forward voltage drop of the diode increases. A drop of 1 V to 1.5 V is typical at full rated current for power diodes.

Shockley diode equation

The Shockley ideal diode equation or the diode law (named after transistor co-inventor William Bradford Shockley, not to be confused with tetrode inventor Walter H. Schottky) gives the I–V characteristic of an ideal diode in either forward or reverse bias (or no bias). The equation is:

where

I is the diode current,

I_S is the reverse bias saturation current,

V_D is the voltage across the diode,

V_T is the thermal voltage, and

n is the emission coefficient, also known as the ideality factor. The emission coefficient n varies from about 1 to 2 depending on the fabrication process and semiconductor material and in many cases is assumed to be approximately equal to 1 (thus the notation n is omitted).

The thermal voltage V_T is approximately 25.85 mV at 300 K, a temperature close to “room temperature” commonly used in device simulation software. At any temperature it is a known constant defined by:

where k is the Boltzmann constant, T is the absolute temperature of the p-n junction, and q is the magnitude of charge on an electron (the elementary charge).
The Shockley ideal diode equation or the diode law is derived with the assumption that the only processes giving rise to the current in the diode are drift (due to electrical field), diffusion, and thermal recombination-generation. It also assumes that the recombination-generation (R-G) current in the depletion region is insignificant. This means that the Shockley equation doesn’t account for the processes involved in reverse breakdown and photon-assisted R-G. Additionally, it doesn’t describe the “leveling off” of the I–V curve at high forward bias due to internal resistance.
Under reverse bias voltages (see Figure 5) the exponential in the diode equation is negligible, and the current is a constant (negative) reverse current value of −I_S. The reverse breakdown region is not modeled by the Shockley diode equation.
For even rather small forward bias voltages (see Figure 5) the exponential is very large because the thermal voltage is very small, so the subtracted ‘1’ in the diode equation is negligible and the forward diode current is often approximated as

The use of the diode equation in circuit problems is illustrated in the article on diode modeling.

Small-signal behaviour

For circuit design, a small-signal model of the diode behavior often proves useful. A specific example of diode modeling is discussed in the article on small-signal circuits.

Types of semiconductor diode


Diode	Zener diode	Schottky diode	Tunnel diode

Light-emitting diode	Photodiode	Varicap	Silicon controlled rectifier

Figure 6: Some diode symbols.

Figure 7: Typical diode packages in same alignment as diode symbol. Thin bar depicts the cathode.

Figure 8: Several types of diodes. The scale is centimeters.

There are several types of junction diodes, which either emphasize a different physical aspect of a diode often by geometric scaling, doping level, choosing the right electrodes, are just an application of a diode in a special circuit, or are really different devices like the Gunn and laser diode and the MOSFET:
Normal (p-n) diodes, which operate as described above, are usually made of doped silicon or, more rarely, germanium. Before the development of modern silicon power rectifier diodes, cuprous oxide and later selenium was used; its low efficiency gave it a much higher forward voltage drop (typically 1.4–1.7 V per “cell”, with multiple cells stacked to increase the peak inverse voltage rating in high voltage rectifiers), and required a large heat sink (often an extension of the diode’s metal substrate), much larger than a silicon diode of the same current ratings would require. The vast majority of all diodes are the p-n diodes found in CMOS integrated circuits, which include two diodes per pin and many other internal diodes.
Avalanche diodes

Diodes that conduct in the reverse direction when the reverse bias voltage exceeds the breakdown voltage. These are electrically very similar to Zener diodes, and are often mistakenly called Zener diodes, but break down by a different mechanism, the avalanche effect. This occurs when the reverse electric field across the p-n junction causes a wave of ionization, reminiscent of an avalanche, leading to a large current. Avalanche diodes are designed to break down at a well-defined reverse voltage without being destroyed. The difference between the avalanche diode (which has a reverse breakdown above about 6.2 V) and the Zener is that the channel length of the former exceeds the “mean free path” of the electrons, so there are collisions between them on the way out. The only practical difference is that the two types have temperature coefficients of opposite polarities.

Cat’s whisker or crystal diodes

These are a type of point-contact diode. The cat’s whisker diode consists of a thin or sharpened metal wire pressed against a semiconducting crystal, typically galena or a piece of coal.^[7] The wire forms the anode and the crystal forms the cathode. Cat’s whisker diodes were also called crystal diodes and found application in crystal radio receivers. Cat’s whisker diodes are generally obsolete, but may be available from a few manufacturers.^{[citation needed]}

Constant current diodes

These are actually a JFET with the gate shorted to the source, and function like a two-terminal current-limiter analog to the Zener diode, which is limiting voltage. They allow a current through them to rise to a certain value, and then level off at a specific value. Also called CLDs, constant-current diodes, diode-connected transistors, or current-regulating diodes.

Esaki or tunnel diodes

These have a region of operation showing negative resistance caused by quantum tunneling, thus allowing amplification of signals and very simple bistable circuits. These diodes are also the type most resistant to nuclear radiation.

Gunn diodes

These are similar to tunnel diodes in that they are made of materials such as GaAs or InP that exhibit a region of negative differential resistance. With appropriate biasing, dipole domains form and travel across the diode, allowing high frequency microwave oscillators to be built.

Light-emitting diodes (LEDs)

In a diode formed from a direct band-gap semiconductor, such as gallium arsenide, carriers that cross the junction emit photons when they recombine with the majority carrier on the other side. Depending on the material, wavelengths (or colors) from the infrared to the near ultraviolet may be produced. The forward potential of these diodes depends on the wavelength of the emitted photons: 1.2 V corresponds to red, 2.4 V to violet. The first LEDs were red and yellow, and higher-frequency diodes have been developed over time. All LEDs produce incoherent, narrow-spectrum light; “white” LEDs are actually combinations of three LEDs of a different color, or a blue LED with a yellow scintillator coating. LEDs can also be used as low-efficiency photodiodes in signal applications. An LED may be paired with a photodiode or phototransistor in the same package, to form an opto-isolator.

Laser diodes

When an LED-like structure is contained in a resonant cavity formed by polishing the parallel end faces, a laser can be formed. Laser diodes are commonly used in optical storage devices and for high speed optical communication.

Peltier diodes

These diodes are used as sensors, heat engines for thermoelectric cooling. Charge carriers absorb and emit their band gap energies as heat.

Photodiodes

All semiconductors are subject to optical charge carrier generation. This is typically an undesired effect, so most semiconductors are packaged in light blocking material. Photodiodes are intended to sense light(photodetector), so they are packaged in materials that allow light to pass, and are usually PIN (the kind of diode most sensitive to light). A photodiode can be used in solar cells, in photometry, or in optical communications. Multiple photodiodes may be packaged in a single device, either as a linear array or as a two-dimensional array. These arrays should not be confused with charge-coupled devices.

Point-contact diodes

These work the same as the junction semiconductor diodes described above, but their construction is simpler. A block of n-type semiconductor is built, and a conducting sharp-point contact made with some group-3 metal is placed in contact with the semiconductor. Some metal migrates into the semiconductor to make a small region of p-type semiconductor near the contact. The long-popular 1N34 germanium version is still used in radio receivers as a detector and occasionally in specialized analog electronics.

PIN diodes

A PIN diode has a central un-doped, or intrinsic, layer, forming a p-type/intrinsic/n-type structure. They are used as radio frequency switches and attenuators. They are also used as large volume ionizing radiation detectors and as photodetectors. PIN diodes are also used in power electronics, as their central layer can withstand high voltages. Furthermore, the PIN structure can be found in many power semiconductor devices, such as IGBTs, power MOSFETs, and thyristors.

Schottky diodes

Schottky diodes are constructed from a metal to semiconductor contact. They have a lower forward voltage drop than p-n junction diodes. Their forward voltage drop at forward currents of about 1 mA is in the range 0.15 V to 0.45 V, which makes them useful in voltage clamping applications and prevention of transistor saturation. They can also be used as low loss rectifiers although their reverse leakage current is generally higher than that of other diodes. Schottky diodes are majority carrier devices and so do not suffer from minority carrier storage problems that slow down many other diodes — so they have a faster “reverse recovery” than p-n junction diodes. They also tend to have much lower junction capacitance than p-n diodes which provides for high switching speeds and their use in high-speed circuitry and RF devices such as switched-mode power supply, mixers and detectors.

Super Barrier Diodes

Super barrier diodes are rectifier diodes that incorporate the low forward voltage drop of the Schottky diode with the surge-handling capability and low reverse leakage current of a normal p-n junction diode.

Gold-doped diodes

As a dopant, gold (or platinum) acts as recombination centers, which help a fast recombination of minority carriers. This allows the diode to operate at signal frequencies, at the expense of a higher forward voltage drop. Gold doped diodes are faster than other p-n diodes (but not as fast as Schottky diodes). They also have less reverse-current leakage than Schottky diodes (but not as good as other p-n diodes).^[8]^[9] A typical example is the 1N914.

Snap-off or Step recovery diodes

The term step recovery relates to the form of the reverse recovery characteristic of these devices. After a forward current has been passing in an SRD and the current is interrupted or reversed, the reverse conduction will cease very abruptly (as in a step waveform). SRDs can therefore provide very fast voltage transitions by the very sudden disappearance of the charge carriers.

Transient voltage suppression diode (TVS)

These are avalanche diodes designed specifically to protect other semiconductor devices from high-voltage transients. Their p-n junctions have a much larger cross-sectional area than those of a normal diode, allowing them to conduct large currents to ground without sustaining damage.

Varicap or varactor diodes

These are used as voltage-controlled capacitors. These are important in PLL (phase-locked loop) and FLL (frequency-locked loop) circuits, allowing tuning circuits, such as those in television receivers, to lock quickly, replacing older designs that took a long time to warm up and lock. A PLL is faster than an FLL, but prone to integer harmonic locking (if one attempts to lock to a broadband signal). They also enabled tunable oscillators in early discrete tuning of radios, where a cheap and stable, but fixed-frequency, crystal oscillator provided the reference frequency for a voltage-controlled oscillator.

Zener diodes

Diodes that can be made to conduct backwards. This effect, called Zener breakdown, occurs at a precisely defined voltage, allowing the diode to be used as a precision voltage reference. In practical voltage reference circuits Zener and switching diodes are connected in series and opposite directions to balance the temperature coefficient to near zero. Some devices labeled as high-voltage Zener diodes are actually avalanche diodes (see above). Two (equivalent) Zeners in series and in reverse order, in the same package, constitute a transient absorber (or Transorb, a registered trademark). The Zener diode is named for Dr. Clarence Melvin Zener of Southern Illinois University, inventor of the device.

Other uses for semiconductor diodes include sensing temperature, and computing analog logarithms (see Operational amplifier applications#Logarithmic).

Numbering and Coding schemes

There are a number of common, standard and manufacturer-driven numbering and coding schemes for diodes; the two most common being the EIA/JEDEC standard and the European Pro Electron standard:

EIA/JEDEC

A standardized 1N-series numbering system was introduced in the US by EIA/JEDEC (Joint Electron Device Engineering Council) about 1960. Among the most popular in this series were: 1N34A/1N270 (Germanium signal), 1N914/1N4148 (Silicon signal), 1N4001-1N4007 (Silicon 1A power rectifier) and 1N54xx (Silicon 3A power rectifier)^[10]^[11]^[12]

Pro Electron

The European Pro Electron coding system for active components was introduced in 1966 and comprises two letters followed by the part code. The first letter represents the semiconductor material used for the component (A = Germanium and B = Silicon) and the second letter represents the general function of the part (for diodes: A = low-power/signal, B = Variable capacitance, X = Multiplier, Y = Rectifier and Z = Voltage reference), for example:

AA-series germanium low-power/signal diodes (e.g.: AA119)
BA-series silicon low-power/signal diodes (e.g.: BAT18 Silicon RF Switching Diode)
BY-series silicon rectifier diodes (e.g.: BY127 1250V, 1A rectifier diode)
BZ-series silicon zener diodes (e.g.: BZY88C4V7 4.7V zener diode)

Other common numbering / coding systems (generally manufacturer-driven) include:

GD-series germanium diodes (ed: GD9) — this is a very old coding system
OA-series germanium diodes (e.g.: OA47) — a coding sequence developed by Mullard, a UK company

As well as these common codes, many manufacturers or organisations have their own systems too — for example:

HP diode 1901-0044 = JEDEC 1N4148
UK military diode CV448 = Mullard type OA81 = GEC type GEX23

Related devices

Rectifier
Transistor
Thyristor or silicon controlled rectifier (SCR)
TRIAC
Diac
Varistor

In optics, an equivalent device for the diode but with laser light would be the Optical isolator, also known as an Optical Diode, that allows light to only pass in one direction. It uses a Faraday rotator as the main component.

Applications

Radio demodulation

The first use for the diode was the demodulation of amplitude modulated (AM) radio broadcasts. The history of this discovery is treated in depth in the radio article. In summary, an AM signal consists of alternating positive and negative peaks of voltage, whose amplitude or “envelope” is proportional to the original audio signal. The diode (originally a crystal diode) rectifies the AM radio frequency signal, leaving an audio signal which is the original audio signal, minus atmospheric noise. The audio is extracted using a simple filter and fed into an audio amplifier or transducer, which generates sound waves.

Power conversion

Rectifiers are constructed from diodes, where they are used to convert alternating current (AC) electricity into direct current (DC). Automotive alternators are a common example, where the diode, which rectifies the AC into DC, provides better performance than the commutator of earlier dynamo. Similarly, diodes are also used in Cockcroft–Walton voltage multipliers to convert AC into higher DC voltages.

Over-voltage protection

Diodes are frequently used to conduct damaging high voltages away from sensitive electronic devices. They are usually reverse-biased (non-conducting) under normal circumstances. When the voltage rises above the normal range, the diodes become forward-biased (conducting). For example, diodes are used in (stepper motor and H-bridge) motor controller and relay circuits to de-energize coils rapidly without the damaging voltage spikes that would otherwise occur. (Any diode used in such an application is called a flyback diode). Many integrated circuits also incorporate diodes on the connection pins to prevent external voltages from damaging their sensitive transistors. Specialized diodes are used to protect from over-voltages at higher power (see Diode types above).

Logic gates

Diodes can be combined with other components to construct AND and OR logic gates. This is referred to as diode logic.

Ionizing radiation detectors

In addition to light, mentioned above, semiconductor diodes are sensitive to more energetic radiation. In electronics, cosmic rays and other sources of ionizing radiation cause noise pulses and single and multiple bit errors. This effect is sometimes exploited by particle detectors to detect radiation. A single particle of radiation, with thousands or millions of electron volts of energy, generates many charge carrier pairs, as its energy is deposited in the semiconductor material. If the depletion layer is large enough to catch the whole shower or to stop a heavy particle, a fairly accurate measurement of the particle’s energy can be made, simply by measuring the charge conducted and without the complexity of a magnetic spectrometer or etc. These semiconductor radiation detectors need efficient and uniform charge collection and low leakage current. They are often cooled by liquid nitrogen. For longer range (about a centimetre) particles they need a very large depletion depth and large area. For short range particles, they need any contact or un-depleted semiconductor on at least one surface to be very thin. The back-bias voltages are near breakdown (around a thousand volts per centimetre). Germanium and silicon are common materials. Some of these detectors sense position as well as energy. They have a finite life, especially when detecting heavy particles, because of radiation damage. Silicon and germanium are quite different in their ability to convert gamma rays to electron showers.
Semiconductor detectors for high energy particles are used in large numbers. Because of energy loss fluctuations, accurate measurement of the energy deposited is of less use.

Temperature measurements

A diode can be used as a temperature measuring device, since the forward voltage drop across the diode depends on temperature, as in a Silicon bandgap temperature sensor. From the Shockley ideal diode equation given above, it appears the voltage has a positive temperature coefficient (at a constant current) but depends on doping concentration and operating temperature (Sze 2007). The temperature coefficient can be negative as in typical thermistors or positive for temperature sense diodes down to about 20 kelvins. Typically, silicon diodes have approximately −2 mV/˚C temperature coefficient at room temperature.

Current steering

Diodes will prevent currents in unintended directions. To supply power to an electrical circuit during a power failure, the circuit can draw current from a battery. An Uninterruptible power supply may use diodes in this way to ensure that current is only drawn from the battery when necessary. Similarly, small boats typically have two circuits each with their own battery/batteries: one used for engine starting; one used for domestics. Normally both are charged from a single alternator, and a heavy duty split charge diode is used to prevent the higher charge battery (typically the engine battery) from discharging through the lower charged battery when the alternator is not running.
Diodes are also used in electronic musical keyboards. To reduce the amount of wiring needed in electronic musical keyboards, these instruments often use keyboard matrix circuits. The keyboard controller scans the rows and columns to determine which note the player has pressed. The problem with matrix circuits is that when several notes are pressed at once, the current can flow backwards through the circuit and trigger "phantom keys" that cause “ghost” notes to play. To avoid triggering unwanted notes, most keyboard matrix circuits have diodes soldered with the switch under each key of the musical keyboard. The same principle is also used for the switch matrix in solid state pinball machines.

Assembly language

2010-05-16T01:29:00.000-07:00

Assembly language

Assembly languages are a type of low-level languages for programming computers, microprocessors, microcontrollers, and other (usually) integrated circuits. They implement a symbolic representation of the numeric machine codes and other constants needed to program a particular CPU architecture. This representation is usually defined by the hardware manufacturer, and is based on abbreviations (called mnemonics) that help the programmer remember individual instructions, registers, etc. An assembly language family is thus specific to a certain physical (or virtual) computer architecture. This is in contrast to most high-level languages, which are (ideally) portable.
A utility program called an assembler is used to translate assembly language statements into the target computer's machine code. The assembler performs a more or less isomorphic translation (a one-to-one mapping) from mnemonic statements into machine instructions and data. This is in contrast with high-level languages, in which a single statement generally results in many machine instructions.
Many sophisticated assemblers offer additional mechanisms to facilitate program development, control the assembly process, and aid debugging. In particular, most modern assemblers include a macro facility (described below), and are called macro assemblers.

Key concepts

Assembler

Typically a modern assembler creates object code by translating assembly instruction mnemonics into opcodes, and by resolving symbolic names for memory locations and other entities.^[1] The use of symbolic references is a key feature of assemblers, saving tedious calculations and manual address updates after program modifications. Most assemblers also include macro facilities for performing textual substitution—e.g., to generate common short sequences of instructions as inline, instead of called subroutines, or even generate entire programs or program suites.
Assemblers are generally simpler to write than compilers for high-level languages, and have been available since the 1950s. Modern assemblers, especially for RISC based architectures, such as MIPS, Sun SPARC, and HP PA-RISC, as well as x86(-64), optimize instruction scheduling to exploit the CPU pipeline efficiently.
There are two types of assemblers based on how many passes through the source are needed to produce the executable program.

One-pass assemblers go through the source code once and assumes that all symbols will be defined before any instruction that references them.
Two-pass assemblers (and multi-pass assemblers) create a table with all unresolved symbols in the first pass, then use the 2nd pass to resolve these addresses. The advantage of a one-pass assembler is speed, which is not as important as it once was with advances in computer speed and capabilities. The advantage of the two-pass assembler is that symbols can be defined anywhere in the program source. As a result, the program can be defined in a more logical and meaningful way. This makes two-pass assembler programs easier to read and maintain.^[2]

More sophisticated high-level assemblers provide language abstractions such as:

Advanced control structures
High-level procedure/function declarations and invocations
High-level abstract data types, including structures/records, unions, classes, and sets
Sophisticated macro processing (although available on ordinary assemblers since late 1960s for IBM/360, amongst other machines)
Object-Oriented features such as encapsulation, polymorphism, inheritance, interfaces

See Language design below for more details.
Note that, in normal professional usage, the term assembler is often used ambiguously: It is frequently used to refer to an assembly language itself, rather than to the assembler utility. Thus: "CP/CMS was written in S/360 assembler" as opposed to "ASM-H was a widely-used S/370 assembler."^{[citation needed]}

Assembly language

A program written in assembly language consists of a series of instructions--mnemonics that correspond to a stream of executable instructions, when translated by an assembler, that can be loaded into memory and executed.
For example, an x86/IA-32 processor can execute the following binary instruction ('MOV') as expressed in machine language (see x86 assembly language):

Hexadecimal: B0 61     (Binary: 10110000 01100001)

The equivalent assembly language representation is easier to remember (example in Intel syntax, more mnemonic):

MOV AL, 61h

This instruction means:

Move (really a copy) the hexadecimal value '61' into the processor register known as "AL". (The h-suffix means hexadecimal or = 97 in decimal)

The mnemonic "mov" represents the opcode 1011 which actually copies the value in the second operand into the register indicated by the first operand. The mnemonic was chosen by the designer of the instruction set to abbreviate "move", making it easier for the programmer to remember. Typical of an assembly language statement, a comma-separated list of arguments or parameters follows the opcode.
In practice many programmers drop the word mnemonic and, technically incorrectly, call "mov" an opcode. When they do this they are referring to the underlying binary code which it represents. To put it another way, a mnemonic such as "mov" is not an opcode, but as it symbolizes an opcode, one might refer to "the opcode mov" for example when one intends to refer to the binary opcode it symbolizes rather than to the symbol -- the mnemonic -- itself. As few modern programmers have need to be mindful of actually what binary patterns are (the opcodes for specific instructions), the distinction has in practice become a bit blurred among programmers but not among processor designers^{[citation needed]}.
Transforming assembly into machine language is accomplished by an assembler, and the (partial) reverse by a disassembler. Unlike high-level languages, there is usually a one-to-one correspondence between simple assembly statements and machine language instructions. However, in some cases, an assembler may provide pseudoinstructions (essentially macros) which expand into several machine language instructions to provide commonly needed functionality. For example, for a machine that lacks a "branch if greater or equal" instruction, an assembler may provide a pseudoinstruction that expands to the machine's "set if less than" and "branch if zero (on the result of the set instruction)". Most full-featured assemblers also provide a rich macro language (discussed below) which is used by vendors and programmers to generate more complex code and data sequences.
Each computer architecture and processor architecture usually has its own machine language. On this level, each instruction is simple enough to be executed using a relatively small number of electronic circuits. Computers differ by the number and type of operations they support. For example, a new 64-bit machine would have different circuitry from a 32-bit machine. They may also have different sizes and numbers of registers, and different representations of data types in storage. While most general-purpose computers are able to carry out essentially the same functionality, the ways they do so differ; the corresponding assembly languages reflect these differences.
Multiple sets of mnemonics or assembly-language syntax may exist for a single instruction set, typically instantiated in different assembler programs. In these cases, the most popular one is usually that supplied by the manufacturer and used in its documentation.

Language design

Basic elements

Any Assembly language consists of 3 types of instruction statements which are used to define the program operations:

Opcode mnemonics
Data sections
Assembly directives

Opcode mnemonics

Instructions (statements) in assembly language are generally very simple, unlike those in high-level languages. Generally, an opcode is a symbolic name for a single executable machine language instruction, and there is at least one opcode mnemonic defined for each machine language instruction. Each instruction typically consists of an operation or opcode plus zero or more operands. Most instructions refer to a single value, or a pair of values. Operands can be either immediate (typically one byte values, coded in the instruction itself) or the addresses of data located elsewhere in storage. This is determined by the underlying processor architecture: the assembler merely reflects how this architecture works.

Data sections

There are instructions used to define data elements to hold data and variables. They define the type of data, the length and the alignment of data. These instructions can also define whether the data is available to outside programs (programs assembled separately) or only to the program in which the data section is defined.

Assembly directives and pseudo-ops

Assembly directives are instructions that are executed by the assembler at assembly time, not by the CPU at run time. They can make the assembly of the program dependent on parameters input by the programmer, so that one program can be assembled different ways, perhaps for different applications. They also can be used to manipulate presentation of the program to make it easier for the programmer to read and maintain.
(For example, pseudo-ops would be used to reserve storage areas and optionally their initial contents.) The names of pseudo-ops often start with a dot to distinguish them from machine instructions.
Some assemblers also support pseudo-instructions, which generate two or more machine instructions.
Symbolic assemblers allow programmers to associate arbitrary names (labels or symbols) with memory locations. Usually, every constant and variable is given a name so instructions can reference those locations by name, thus promoting self-documenting code. In executable code, the name of each subroutine is associated with its entry point, so any calls to a subroutine can use its name. Inside subroutines, GOTO destinations are given labels. Some assemblers support local symbols which are lexically distinct from normal symbols (e.g., the use of "10$" as a GOTO destination).
Most assemblers provide flexible symbol management, allowing programmers to manage different namespaces, automatically calculate offsets within data structures, and assign labels that refer to literal values or the result of simple computations performed by the assembler. Labels can also be used to initialize constants and variables with relocatable addresses.
Assembly languages, like most other computer languages, allow comments to be added to assembly source code that are ignored by the assembler. Good use of comments is even more important with assembly code than with higher-level languages, as the meaning and purpose of a sequence of instructions is harder to decipher from the code itself.
Wise use of these facilities can greatly simplify the problems of coding and maintaining low-level code. Raw assembly source code as generated by compilers or disassemblers—code without any comments, meaningful symbols, or data definitions—is quite difficult to read when changes must be made.

Macros

Many assemblers support predefined macros, and others support programmer-defined (and repeatedly redefinable) macros involving sequences of text lines that variables and constants are embedded in. This sequence of text lines may include a sequence of instructions, or a sequence of data storage pseudo-ops. Once a macro has been defined using the appropriate pseudo-op, its name may be used in place of a mnemonic. When the assembler processes such a statement, it replaces the statement with the text lines associated with that macro, then processes them just as though they had appeared in the source code file all along (including, in better assemblers, expansion of any macros appearing in the replacement text).
Since macros can have 'short' names but expand to several or indeed many lines of code, they can be used to make assembly language programs appear to be much shorter (require less lines of source code from the application programmer, as with a higher level language). They can also be used to add higher levels of structure to assembly programs, optionally introduce embedded de-bugging code via parameters and other similar features.
Many assemblers have built-in (or predefined) macros for system calls and other special code sequences, such as the generation and storage of data realized through advanced bitwise and boolean operations used in gaming, software security, data management, and cryptography.
Macro assemblers often allow macros to take parameters. Some assemblers include quite sophisticated macro languages, incorporating such high-level language elements as optional parameters, symbolic variables, conditionals, string manipulation, and arithmetic operations, all usable during the execution of a given macro, and allowing macros to save context or exchange information. Thus a macro might generate a large number of assembly language instructions or data definitions, based on the macro arguments. This could be used to generate record-style data structures or "unrolled" loops, for example, or could generate entire algorithms based on complex parameters. An organization using assembly language that has been heavily extended using such a macro suite can be considered to be working in a higher-level language, since such programmers are not working with a computer's lowest-level conceptual elements.
Macros were used to customize large scale software systems for specific customers in the mainframe era and were also used by customer personnel to satisfy their employers' needs by making specific versions of manufacturer operating systems; this was done, for example, by systems programmers working with IBM's Conversational Monitor System/Virtual Machine (CMS/VM) and with IBM's "real time transaction processing" add-ons, CICS, Customer Information Control System, and ACP/TPF, the airline/financial system that began in the 1970s and still runs many large Global Distribution Systems (GDS) and credit card systems today.
It was also possible to use solely the macro processing capabilities of an assembler to generate code written in completely different languages, for example, to generate a version of a program in Cobol using a pure macro assembler program containing lines of Cobol code inside assembly time operators instructing the assembler to generate arbitrary code.
This was because, as was realized in the 1970s, the concept of "macro processing" is independent of the concept of "assembly", the former being in modern terms more word processing, text processing, than generating object code. The concept of macro processing in fact appeared in and appears in the C programming language, which supports "preprocessor instructions" to set variables, and make conditional tests on their values. Note that unlike certain previous macro processors inside assemblers, the C preprocessor was not Turing-complete because it lacked the ability to either loop or "go to", the latter allowing the programmer to loop.
Despite the power of macro processing, it fell into disuse in high level languages while remaining a perennial for assemblers.
This was because many programmers were rather confused by macro parameter substitution and did not disambiguate macro processing from assembly and execution^{[dubious – discuss]}.
Macro parameter substitution is strictly by name: at macro processing time, the value of a parameter is textually substituted for its name. The most famous class of bugs resulting was the use of a parameter that itself was an expression and not a simple name when the macro writer expected a name. In the macro: foo: macro a load a*b the intention was that the caller would provide the name of a variable, and the "global" variable or constant b would be used to multiply "a". If foo is called with the parameter a-c, an unexpected macro expansion occurs.
To avoid this, users of macro processors learned to religiously parenthesize formal parameters inside macro definitions, and callers had to do the same to their "actual" parameters^{[citation needed]}.
PL/I and C feature macros, but this facility was underused or dangerous when used^{[citation needed]} because they can only manipulate text. On the other hand, homoiconic languages, such as Lisp, Prolog, and Forth, retain the power of assembly language macros because they are able to manipulate their own code as data.

Support for structured programming

Some assemblers have incorporated structured programming elements to encode execution flow. The earliest example of this approach was in the Concept-14 macro set, originally proposed by Dr. H.D. Mills (March, 1970), and implemented by Marvin Kessler at IBM's Federal Systems Division, which extended the S/360 macro assembler with IF/ELSE/ENDIF and similar control flow blocks.^[3] This was a way to reduce or eliminate the use of GOTO operations in assembly code, one of the main factors causing spaghetti code in assembly language. This approach was widely accepted in the early 80s (the latter days of large-scale assembly language use).
A curious design was A-natural, a "stream-oriented" assembler for 8080/Z80 processors^{[citation needed]} from Whitesmiths Ltd. (developers of the Unix-like Idris operating system, and what was reported to be the first commercial C compiler). The language was classified as an assembler, because it worked with raw machine elements such as opcodes, registers, and memory references; but it incorporated an expression syntax to indicate execution order. Parentheses and other special symbols, along with block-oriented structured programming constructs, controlled the sequence of the generated instructions. A-natural was built as the object language of a C compiler, rather than for hand-coding, but its logical syntax won some fans.
There has been little apparent demand for more sophisticated assemblers since the decline of large-scale assembly language development.^[4] In spite of that, they are still being developed and applied in cases where resource constraints or peculiarities in the target system's architecture prevent the effective use of higher-level languages.^[5]

Use of assembly language

Historical perspective

Assembly languages were first developed in the 1950s, when they were referred to as second generation programming languages. They eliminated much of the error-prone and time-consuming first-generation programming needed with the earliest computers, freeing the programmer from tedium such as remembering numeric codes and calculating addresses. They were once widely used for all sorts of programming. However, by the 1980s (1990s on small computers), their use had largely been supplanted by high-level languages^{[citation needed]}, in the search for improved programming productivity. Today, although assembly language is almost always handled and generated by compilers, it is still used for direct hardware manipulation, access to specialized processor instructions, or to address critical performance issues. Typical uses are device drivers, low-level embedded systems, and real-time systems.
Historically, a large number of programs have been written entirely in assembly language. Operating systems were almost exclusively written in assembly language until the widespread acceptance of C in the 1970s and early 1980s. Many commercial applications were written in assembly language as well, including a large amount of the IBM mainframe software written by large corporations. COBOL and FORTRAN eventually displaced much of this work, although a number of large organizations retained assembly-language application infrastructures well into the 90s.
Most early microcomputers relied on hand-coded assembly language, including most operating systems and large applications. This was because these systems had severe resource constraints, imposed idiosyncratic memory and display architectures, and provided limited, buggy system services. Perhaps more important was the lack of first-class high-level language compilers suitable for microcomputer use. A psychological factor may have also played a role: the first generation of microcomputer programmers retained a hobbyist, "wires and pliers" attitude.
In a more commercial context, the biggest reasons for using assembly language were minimal bloat (size), minimal overhead, greater speed, and reliability.
Typical examples of large assembly language programs from this time are the MS-DOS operating system, the early IBM PC spreadsheet program Lotus 1-2-3, and almost all popular games for the Atari 800 family of home computers. Even into the 1990s, most console video games were written in assembly, including most games for the Mega Drive/Genesis and the Super Nintendo Entertainment System^{[citation needed]}. According to some industry insiders, the assembly language was the best computer language to use to get the best performance out of the Sega Saturn, a console that was notoriously challenging to develop and program games for ^[6]. The popular arcade game NBA Jam (1993) is another example. On the Commodore 64, Amiga, Atari ST, as well as ZX Spectrum home computers, assembler has long been the primary development language. This was in large part due to the fact that BASIC dialects on these systems offered insufficient execution speed, as well as insufficient facilities to take full advantage of the available hardware on these systems. Some systems, most notably Amiga, even have IDEs with highly advanced debugging and macro facilities, such as the freeware ASM-One assembler, comparable to that of Microsoft Visual Studio facilities (ASM-One predates Microsoft Visual Studio).
The Assembler for the VIC-20 was written by Don French and published by French Silk. At 1639 bytes in length, its author believes it is the smallest symbolic assembler ever written. The assembler supported the usual symbolic addressing and the definition of character strings or hex strings. It also allowed address expressions which could be combined with addition, subtraction, multiplication, division, logical AND, logical OR, and exponentiation operators.^[7]

Current usage

There have always been debates over the usefulness and performance of assembly language relative to high-level languages. Assembly language has specific niche uses where it is important; see below. But in general, modern optimizing compilers are claimed^{[citation needed]} to render high-level languages into code that can run as fast as hand-written assembly, despite the counter-examples that can be found ^[8]^[9]^[10]. The complexity of modern processors and memory sub-system makes effective optimization increasingly difficult for compilers, as well as assembler programmers ^[11]^[12]. Moreover, and to the dismay of efficiency lovers, increasing processor performance has meant that most CPUs sit idle most of the time, with delays caused by predictable bottlenecks such as I/O operations and paging. This has made raw code execution speed a non-issue for many programmers.
There are some situations in which practitioners might choose to use assembly language, such as when:

a stand-alone binary executable is required, i.e. one that must execute without recourse to the run-time components or libraries associated with a high-level language; this is perhaps the most common situation. These are embedded programs that store only a small amount of memory and the device is intended to do single purpose tasks. Such examples consist of telephones, automobile fuel and ignition systems, air-conditioning control systems, security systems, and sensors.
interacting directly with the hardware, for example in device drivers and interrupt handlers.
using processor-specific instructions not exploited by or available to the compiler. A common example is the bitwise rotation instruction at the core of many encryption algorithms.
creating vectorized functions for programs in higher-level languages such as C. In the higher-level language this is sometimes aided by compiler intrinsic functions which map directly to SIMD mnemonics, but nevertheless result in a one-to-one assembly conversion specific for the given vector processor.
extreme optimization is required, e.g., in an inner loop in a processor-intensive algorithm. Game programmers take advantage of the capabilities of hardware features in systems, enabling the games to run faster.
a system with severe resource constraints (e.g., an embedded system) must be hand-coded to maximize the use of limited resources; but this is becoming less common as processor price decreases and performance improves.
no high-level language exists, on a new or specialized processor, for example.
writing real-time programs that need precise timing and responses, such as simulations, flight navigation systems, and medical equipment. For example, in a fly-by-wire system, telemetry must be interpreted and acted upon within strict time constraints. Such systems must eliminate sources of unpredictable delays, which may be created by (some) interpreted languages, automatic garbage collection, paging operations, or preemptive multitasking. However, some higher-level languages incorporate run-time components and operating system interfaces that can introduce such delays. Choosing assembly or lower-level languages for such systems gives the programmer greater visibility and control over processing details.
complete control over the environment is required, in extremely high security situations where nothing can be taken for granted.
writing computer viruses, bootloaders, certain device drivers, or other items very close to the hardware or low-level operating system.
writing instruction set simulators for monitoring, tracing and debugging where additional overhead is kept to a minimum
reverse-engineering existing binaries that may or may not have originally been written in a high-level language, for example when cracking copy protection of proprietary software.
reverse engineering and modifying video games (also known as ROM Hacking), which is possible with a range of techniques. The most widely employed is altering the program code at the assembly language level.
writing self modifying code, to which assembly language lends itself well.
writing games and other software for graphing calculators.^[13]
writing compiler software that generates assembly code, and the writers should therefore be expert assembly language programmers themselves.
writing cryptographic algorithms that must always take strictly the same time to execute, preventing timing attacks.

Nevertheless, assembly language is still taught in most Computer Science and Electronic Engineering programs. Although few programmers today regularly work with assembly language as a tool, the underlying concepts remain very important. Such fundamental topics as binary arithmetic, memory allocation, stack processing, character set encoding, interrupt processing, and compiler design would be hard to study in detail without a grasp of how a computer operates at the hardware level. Since a computer's behavior is fundamentally defined by its instruction set, the logical way to learn such concepts is to study an assembly language. Most modern computers have similar instruction sets. Therefore, studying a single assembly language is sufficient to learn: i) The basic concepts; ii) To recognize situations where the use of assembly language might be appropriate; and iii) To see how efficient executable code can be created from high-level languages. ^[14]

Typical applications

Hard-coded assembly language is typically used in a system's boot ROM (BIOS on IBM-compatible PC systems). This low-level code is used, among other things, to initialize and test the system hardware prior to booting the OS, and is stored in ROM. Once a certain level of hardware initialization has taken place, execution transfers to other code, typically written in higher level languages; but the code running immediately after power is applied is usually written in assembly language. The same is true of most boot loaders.
Many compilers render high-level languages into assembly first before fully compiling, allowing the assembly code to be viewed for debugging and optimization purposes. Relatively low-level languages, such as C, often provide special syntax to embed assembly language directly in the source code. Programs using such facilities, such as the Linux kernel, can then construct abstractions utilizing different assembly language on each hardware platform. The system's portable code can then utilize these processor-specific components through a uniform interface.
Assembly language is also valuable in reverse engineering, since many programs are distributed only in machine code form, and machine code is usually easy to translate into assembly language and carefully examine in this form, but very difficult to translate into a higher-level language. Tools such as the Interactive Disassembler make extensive use of disassembly for such a purpose.
A particular niche that makes use of assembly language is the demoscene. Certain competitions require the contestants to restrict their creations to a very small size (e.g. 256B, 1KB, 4KB or 64 KB), and assembly language is the language of choice to achieve this goal.^[15] When resources, particularly CPU-processing constrained systems, like the earlier Amiga models, and the Commodore 64, are a concern, assembler coding is a must: optimized assembler code is written "by hand" and instructions are sequenced manually by the coders in an attempt to minimize the number of CPU cycles used; the CPU constraints are so great that every CPU cycle counts. However, using such techniques has enabled systems like the Commodore 64 to produce real-time 3D graphics with advanced effects, a feat which might be considered unlikely or even impossible for a system with a 0.99MHz processor.^{[citation needed]}

Related terminology

Assembly language or assembler language is commonly called assembly, assembler, ASM, or symbolic machine code. A generation of IBM mainframe programmers called it BAL for Basic Assembly Language.

Note: Calling the language assembler is of course potentially confusing and ambiguous, since this is also the name of the utility program that translates assembly language statements into machine code. Some may regard this as imprecision or error. However, this usage has been common among professionals and in the literature for decades.^[16] Similarly, some early computers called their assembler its assembly program.^[17])

The computational step where an assembler is run, including all macro processing, is known as assembly time.
The use of the word assembly dates from the early years of computers (cf. short code, speedcode).
A cross assembler (see cross compiler) is functionally just an assembler. This term is used to stress that the assembler is run on a different computer than the target system, the system on which the resulting code is run. Because nowadays assemblers are written portably in a high level language like C, this is largely irrelevant. Cross assembling may be necessary if the target system lacks the capacity to run an assembler itself. This is typically the case for small embedded systems. The most important distinguishing feature of a cross assembler is that it provides for or interfaces to facilities to transport the code to the target processor, e.g. to reside in flash or EPROM. It generates a binary image, or Intel Hex file rather than an object file.
An assembler directive is a command given to an assembler. These directives may do anything from telling the assembler to include other source files, to telling it to allocate memory for constant data.

List of assemblers for different computer architectures

The following page has a list of different assemblers for the different computer architectures, along with any associated information for that specific assembler:

List of assemblers

Further details

For any given personal computer, mainframe, embedded system, and game console, both past and present, at least one--possibly dozens--of assemblers have been written. For some examples, see the list of assemblers.
On Unix systems, the assembler is traditionally called as, although it is not a single body of code, being typically written anew for each port. A number of Unix variants use GAS.
Within processor groups, each assembler has its own dialect. Sometimes, some assemblers can read another assembler's dialect, for example, TASM can read old MASM code, but not the reverse. FASM and NASM have similar syntax, but each support different macros that could make them difficult to translate to each other. The basics are all the same, but the advanced features will differ.^[18]
Also, assembly can sometimes be portable across different operating systems on the same type of CPU. Calling conventions between operating systems often differ slightly or not at all, and with care it is possible to gain some portability in assembly language, usually by linking with a C library that does not change between operating systems. An instruction set simulator (which would ideally be written in an assembler language) can, in theory, process the object code/ binary of any assembler to achieve portability even across platforms (with an overhead no greater than a typical bytecode interpreter). This is essentially what microcode achieves when a hardware platform changes internally.
For example, many things in libc depend on the preprocessor to do OS-specific, C-specific things to the program before compiling. In fact, some functions and symbols are not even guaranteed to exist outside of the preprocessor. Worse, the size and field order of structs, as well as the size of certain typedefs such as off_t, are entirely unavailable in assembly language without help from a configure script, and differ even between versions of Linux, making it impossible to portably call functions in libc other than ones that only take simple integers and pointers as parameters. To address this issue, FASMLIB project provides a portable assembly library for Win32 and Linux platforms, but it is yet very incomplete.^[19]
Some higher level computer languages, such as C and Borland Pascal, support inline assembly where relatively brief sections of assembly code can be embedded into the high level language code. The Forth programming language commonly contains an assembler used in CODE words.
Many people use an emulator to debug assembly-language programs.

Example listing of assembly language source code

Address	Label	Instruction (AT&T syntax)	Object code^[20]
		`.begin`
		`.org 2048`
	`a_start`	`.equ 3000`
`2048`		`ld length,%`
`2064`		`be done`	`00000010 10000000 00000000 00000110`
`2068`		`addcc %r1,-4,%r1`	`10000010 10000000 01111111 11111100`
`2072`		`addcc %r1,%r2,%r4`	`10001000 10000000 01000000 00000010`
`2076`		`ld %r4,%r5`	`11001010 00000001 00000000 00000000`
`2080`		`ba loop`	`00010000 10111111 11111111 11111011`
`2084`		`addcc %r3,%r5,%r3`	`10000110 10000000 11000000 00000101`
`2088`	`done:`	`jmpl %r15+4,%r0`	`10000001 11000011 11100000 00000100`
`2092`	`length:`	`20`	`00000000 00000000 00000000 00010100`
`2096`	`address:`	`a_start`	`00000000 00000000 00001011 10111000`
		`.org a_start`
`3000`	`a:`

Example of a selection of instructions (for a virtual computer^[21]) with the corresponding address in memory where each instruction will be placed. These addresses are not static, see memory management. Accompanying each instruction is the generated (by the assembler) object code that coincides with the virtual computer's architecture (or ISA).

Microcontroller

2010-05-16T01:05:00.000-07:00

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Please help improve this article by adding reliable references. Unsourced material may be challenged and removed. (November 2009)
The integrated circuit from an Intel 8742, an 8-bit microcontroller that includes a CPU running at 12 MHz, 128 bytes of RAM, 2048 bytes of EPROM, and I/O in the same chip.

A microcontroller is a small computer on a single integrated circuit containing a processor core, memory, and programmable input/output peripherals. Program memory in the form of NOR flash or OTP ROM is also often included on chip, as well as a typically small amount of RAM. Microcontrollers are designed for embedded applications, in contrast to the microprocessors used in personal computers or other general purpose applications.

Microcontrollers are used in automatically controlled products and devices, such as automobile engine control systems, implantable medical devices, remote controls, office machines, appliances, power tools, and toys. By reducing the size and cost compared to a design that uses a separate microprocessor, memory, and input/output devices, microcontrollers make it economical to digitally control even more devices and processes. Mixed signal microcontrollers are common, integrating analog components needed to control non-digital electronic systems.

Some microcontrollers may use four-bit words and operate at clock rate frequencies as low as 4 kHz, for low power consumption (milliwatts or microwatts). They will generally have the ability to retain functionality while waiting for an event such as a button press or other interrupt; power consumption while sleeping (CPU clock and most peripherals off) may be just nanowatts, making many of them well suited for long lasting battery applications. Other microcontrollers may serve performance-critical roles, where they may need to act more like a digital signal processor (DSP), with higher clock speeds and power consumption.

Contents
[hide]

* 1 Embedded design
o 1.1 Interrupts
o 1.2 Programs
o 1.3 Other microcontroller features
* 2 Higher integration
* 3 Volumes
* 4 Programming environments
* 5 Types of microcontrollers
* 6 Interrupt latency
* 7 History
* 8 Microcontroller embedded memory technology
o 8.1 Data
o 8.2 Firmware
* 9 See also
* 10 Notes
* 11 External links

Embedded design

A microcontroller can be considered a self-contained system with a processor, memory and peripherals and can be used as an embedded system. [1] The majority of microcontrollers in use today are embedded in other machinery, such as automobiles, telephones, appliances, and peripherals for computer systems. These are called embedded systems. While some embedded systems are very sophisticated, many have minimal requirements for memory and program length, with no operating system, and low software complexity. Typical input and output devices include switches, relays, solenoids, LEDs, small or custom LCD displays, radio frequency devices, and sensors for data such as temperature, humidity, light level etc. Embedded systems usually have no keyboard, screen, disks, printers, or other recognizable I/O devices of a personal computer, and may lack human interaction devices of any kind.
Interrupts

Microcontrollers must provide real time (predictable, though not necessarily fast) response to events in the embedded system they are controlling. When certain events occur, an interrupt system can signal the processor to suspend processing the current instruction sequence and to begin an interrupt service routine (ISR, or "interrupt handler"). The ISR will perform any processing required based on the source of the interrupt before returning to the original instruction sequence. Possible interrupt sources are device dependent, and often include events such as an internal timer overflow, completing an analog to digital conversion, a logic level change on an input such as from a button being pressed, and data received on a communication link. Where power consumption is important as in battery operated devices, interrupts may also wake a microcontroller from a low power sleep state where the processor is halted until required to do something by a peripheral event.
Programs

Microcontroller programs must fit in the available on-chip program memory, since it would be costly to provide a system with external, expandable, memory. Compilers and assemblers are used to turn high-level language and assembler language codes into a compact machine code for storage in the microcontroller's memory. Depending on the device, the program memory may be permanent, read-only memory that can only be programmed at the factory, or program memory may be field-alterable flash or erasable read-only memory.
Other microcontroller features

Microcontrollers usually contain from several to dozens of general purpose input/output pins (GPIO). GPIO pins are software configurable to either an input or an output state. When GPIO pins are configured to an input state, they are often used to read sensors or external signals. Configured to the output state, GPIO pins can drive external devices such as LED's or motors.

Many embedded systems need to read sensors that produce analog signals. This is the purpose of the analog-to-digital converter (ADC). Since processors are built to interpret and process digital data, i.e. 1s and 0s, they won't be able to do anything with the analog signals that may be sent to it by a device. So the analog to digital converter is used to convert the incoming data into a form that the processor can recognize. A less common feature on some microcontrollers is a digital-to-analog converter (DAC) that allows the processor to output analog signals or voltage levels.

In addition to the converters, many embedded microprocessors include a variety of timers as well. One of the most common types of timers is the Programmable Interval Timer (PIT). A PIT just counts down from some value to zero. Once it reaches zero, it sends an interrupt to the processor indicating that it has finished counting. This is useful for devices such as thermostats, which periodically test the temperature around them to see if they need to turn the air conditioner on, the heater on, etc.

Time Processing Unit (TPU) is a sophisticated timer. In addition to counting down, the TPU can detect input events, generate output events, and perform other useful operations.

A dedicated Pulse Width Modulation (PWM) block makes it possible for the CPU to control power converters, resistive loads, motors, etc., without using lots of CPU resources in tight timer loops.

Universal Asynchronous Receiver/Transmitter (UART) block makes it possible to receive and transmit data over a serial line with very little load on the CPU. Dedicated on-chip hardware also often includes capabilities to communicate with other devices (chips) in digital formats such as I2C and Serial Peripheral Interface (SPI).
Higher integration

In contrast to general-purpose CPUs, micro-controllers may not implement an external address or data bus as they integrate RAM and non-volatile memory on the same chip as the CPU. Using fewer pins, the chip can be placed in a much smaller, cheaper package.

Integrating the memory and other peripherals on a single chip and testing them as a unit increases the cost of that chip, but often results in decreased net cost of the embedded system as a whole. Even if the cost of a CPU that has integrated peripherals is slightly more than the cost of a CPU and external peripherals, having fewer chips typically allows a smaller and cheaper circuit board, and reduces the labor required to assemble and test the circuit board.

A micro-controller is a single integrated circuit, commonly with the following features:

* central processing unit - ranging from small and simple 4-bit processors to complex 32- or 64-bit processors
* discrete input and output bits, allowing control or detection of the logic state of an individual package pin
* serial input/output such as serial ports (UARTs)
* other serial communications interfaces like I²C, Serial Peripheral Interface and Controller Area Network for system interconnect
* peripherals such as timers, event counters, PWM generators, and watchdog
* volatile memory (RAM) for data storage
* ROM, EPROM, EEPROM or Flash memory for program and operating parameter storage
* clock generator - often an oscillator for a quartz timing crystal, resonator or RC circuit
* many include analog-to-digital converters
* in-circuit programming and debugging support

This integration drastically reduces the number of chips and the amount of wiring and circuit board space that would be needed to produce equivalent systems using separate chips. Furthermore, and on low pin count devices in particular, each pin may interface to several internal peripherals, with the pin function selected by software. This allows a part to be used in a wider variety of applications than if pins had dedicated functions. Micro-controllers have proved to be highly popular in embedded systems since their introduction in the 1970s.

Some microcontrollers use a Harvard architecture: separate memory buses for instructions and data, allowing accesses to take place concurrently. Where a Harvard architecture is used, instruction words for the processor may be a different bit size than the length of internal memory and registers; for example: 12-bit instructions used with 8-bit data registers.

The decision of which peripheral to integrate is often difficult. The microcontroller vendors often trade operating frequencies and system design flexibility against time-to-market requirements from their customers and overall lower system cost. Manufacturers have to balance the need to minimize the chip size against additional functionality.

Microcontroller architectures vary widely. Some designs include general-purpose microprocessor cores, with one or more ROM, RAM, or I/O functions integrated onto the package. Other designs are purpose built for control applications. A micro-controller instruction set usually has many instructions intended for bit-wise operations to make control programs more compact.[2] For example, a general purpose processor might require several instructions to test a bit in a register and branch if the bit is set, where a micro-controller could have a single instruction to provide that commonly-required function.

Microcontrollers typically do not have a math coprocessor, so floating point arithmetic is performed by software.
Volumes

About 55% of all CPUs sold in the world are 8-bit microcontrollers and microprocessors. According to Semico, over four billion 8-bit microcontrollers were sold in 2006.[3]

A typical home in a developed country is likely to have only four general-purpose microprocessors but around three dozen microcontrollers. A typical mid-range automobile has as many as 30 or more microcontrollers. They can also be found in many electrical device such as washing machines, microwave ovens, and telephones.
A PIC 18F8720 microcontroller in an 80-pin TQFP package.

Manufacturers have often produced special versions of their microcontrollers in order to help the hardware and software development of the target system. Originally these included EPROM versions that have a "window" on the top of the device through which program memory can be erased by ultraviolet light, ready for reprogramming after a programming ("burn") and test cycle. Since 1998, EPROM versions are rare and have been replaced by EEPROM and flash, which are easier to use (can be erased electronically) and cheaper to manufacture.

Other versions may be available where the ROM is accessed as an external device rather than as internal memory, however these are becoming increasingly rare due to the widespread availability of cheap microcontroller programmers.

The use of field-programmable devices on a microcontroller may allow field update of the firmware or permit late factory revisions to products that have been assembled but not yet shipped. Programmable memory also reduces the lead time required for deployment of a new product.

Where hundreds of thousands of identical devices are required, using parts programmed at the time of manufacture can be an economical option. These 'mask programmed' parts have the program laid down in the same way as the logic of the chip, at the same time.
[edit] Programming environments

Microcontrollers were originally programmed only in assembly language, but various high-level programming languages are now also in common use to target microcontrollers. These languages are either designed specially for the purpose, or versions of general purpose languages such as the C programming language. Compilers for general purpose languages will typically have some restrictions as well as enhancements to better support the unique characteristics of microcontrollers. Some microcontrollers have environments to aid developing certain types of applications. Microcontroller vendors often make tools freely available to make it easier to adopt their hardware.

Many microcontrollers are so quirky that they effectively require their own non-standard dialects of C, such as SDCC for the 8051, which prevent using standard tools (such as code libraries or static analysis tools) even for code unrelated to hardware features. Interpreters are often used to hide such low level quirks.

Interpreter firmware is also available for some microcontrollers. For example, BASIC on the early microcontrollers Intel 8052[4]; BASIC and FORTH on the Zilog Z8[5] as well as some modern devices. Typically these interpreters support interactive programming.

Simulators are available for some microcontrollers, such as in Microchip's MPLAB environment. These allow a developer to analyze what the behavior of the microcontroller and their program should be if they were using the actual part. A simulator will show the internal processor state and also that of the outputs, as well as allowing input signals to be generated. While on the one hand most simulators will be limited from being unable to simulate much other hardware in a system, they can exercise conditions that may otherwise be hard to reproduce at will in the physical implementation, and can be the quickest way to debug and analyze problems.

Recent microcontrollers are often integrated with on-chip debug circuitry that when accessed by an in-circuit emulator via JTAG, allow debugging of the firmware with a debugger.
Types of microcontrollers
Wiki letter w.svg This section requires expansion.
See also: List of common microcontrollers

As of 2008 there are several dozen microcontroller architectures and vendors including:

* 68HC11
* 8051
* ARM processors (from many vendors) using ARM7 or Cortex-M3 cores are generally microcontrollers
* STMicroelectronics STM8S (8-bit), ST10 (16-bit) and STM32 (32-bit)
* Atmel AVR (8-bit), AVR32 (32-bit), and AT91SAM
* Freescale ColdFire (32-bit) and S08 (8-bit)
* Hitachi H8, Hitachi SuperH
* Hyperstone E1/E2 (32-bit, First full integration of RISC and DSP on one processor core [1996] [1])
* MIPS (32-bit PIC32)
* NEC V850
* PIC (8-bit PIC16, PIC18, 16-bit dsPIC33 / PIC24)
* PowerPC ISE
* PSoC (Programmable System-on-Chip)
* Rabbit 2000
* Texas Instruments MSP430 (16-bit), C2000 (32-bit), and Stellaris (32-bit)
* Toshiba TLCS-870
* Zilog eZ8, eZ80

and many others, some of which are used in very narrow range of applications or are more like applications processors than microcontrollers. The microcontroller market is extremely fragmented, with numerous vendors, technologies, and markets. Note that many vendors sell (or have sold) multiple architectures.
Interrupt latency

In contrast to general-purpose computers, microcontrollers used in embedded systems often seek to optimize interrupt latency over instruction throughput. Issues include both reducing the latency, and making it be more predictable (to support real-time control).

When an electronic device causes an interrupt, the intermediate results (registers) have to be saved before the software responsible for handling the interrupt can run. They must also be restored after that software is finished. If there are more registers, this saving and restoring process takes more time, increasing the latency. Ways to reduce such context/restore latency include having relatively few registers in their central processing units (undesirable because it slows down most non-interrupt processing substantially), or at least not having hardware save them all (hoping that the software doesn't then need to compensate by saving the rest "manually"). Another technique involves spending silicon gates on "shadow registers": one or more duplicate registers used only by the interrupt software, perhaps supporting a dedicated stack.

Other factors affecting interrupt latency include:

* Cycles needed to complete current CPU activities. To minimize those costs, microcontrollers tend to have short pipelines (often three instructions or less), small write buffers, and ensure that longer instructions are continuable or restartable. RISC design principles ensure that most instructions take the same number of cycles, helping avoid the need for most such continuation/restart logic.
* The length of any critical section that needs to be interrupted. Entry to a critical section restricts concurrent data structure access. When a data structure must be accessed by an interrupt handler, the critical section must block that interrupt. Accordingly, interrupt latency is increased by however long that interrupt is blocked. When there are hard external constraints on system latency, developers often need tools to measure interrupt latencies and track down which critical sections cause slowdowns.
o One common technique just blocks all interrupts for the duration of the critical section. This is easy to implement, but sometimes critical sections get uncomfortably long.
o A more complex technique just blocks the interrupts that may trigger access to that data structure. This often based on interrupt priorities, which tend to not correspond well to the relevant system data structures. Accordingly, this technique is used mostly in very constrained environments.
o Processors may have hardware support for some critical sections. Examples include supporting atomic access to bits or bytes within a word, or other atomic access primitives like the LDREX/STREX exclusive access primitives introduced in the ARMv6 architecture.
* Interrupt nesting. Some microcontrollers allow higher priority interrupts to interrupt lower priority ones. This allows software to manage latency by giving time-critical interrupts higher priority (and thus lower and more predictable latency) than less-critical ones.
* Trigger rate. When interrupts occur back-to-back, microcontrollers may avoid an extra context save/restore cycle by a form of tail call optimization.

Lower end microcontrollers tend to support fewer interrupt latency controls than higher end ones.
History
Wiki letter w.svg This section requires expansion.

The first single-chip microprocessor was the 4-bit Intel 4004 released in 1971, headed by Intels lead research scientist Hunter H. Hetfeld. With the Intel 8008 and more capable microprocessors available over the next several years.

These however all required external chip(s) to implement a working system, raising total system cost, and making it impossible to economically computerise appliances.

The first computer system on a chip optimised for control applications - microcontroller was the Intel 8048 released in 1975[citation needed], with both RAM and ROM on the same chip. This chip would find its way into over one billion PC keyboards, and other numerous applications. At this time Intels President, Luke J. Valenter, stated that the (Microcontroller) was one of the most successful in the companies history, and expanded the division's budget over 25%.

Most microcontrollers at this time had two variants. One had an erasable EEPROM program memory, which was significantly more expensive than the PROM variant which was only programmable once.

In 1993, the introduction of EEPROM memory allowed microcontrollers (beginning with the Microchip PIC16x84) [2][citation needed]) to be electrically erased quickly without an expensive package as required for EPROM, allowing both rapid prototyping, and In System Programming.

The same year, Atmel introduced the first microcontroller using Flash memory. [6].

Other companies rapidly followed suit, with both memory types.

Cost has plummeted over time, with the cheapest 8-bit microcontrollers being available for under $0.25 in quantity (thousands) in 2009, and some 32-bit microcontrollers around $1 for similar quantities.

Nowadays microcontrollers are low cost and readily available for hobbyists, with large online communities around certain processors.

In the future, MRAM could potentially be used in microcontrollers as it has infinite endurance and its incremental semiconductor wafer process cost is relatively low.
Microcontroller embedded memory technology

Since the emergence of microcontrollers, many different memory technologies have been used. Almost all microcontrollers have at least two different kinds of memory, a non-volatile memory for storing firmware and a read-write memory for temporary data.
Data

From the earliest microcontrollers to today, six-transistor SRAM is almost always used as the read/write working memory, with a few more transistors per bit used in the register file. MRAM could potentially replace it as it is 4-10 times denser which would make it more cost effective.

In addition to the SRAM, some microcontrollers also have internal EEPROM for data storage; and even ones that don't have any (or don't have enough) are often connected to external serial EEPROM chip (such as the BASIC Stamp) or external serial flash memory chip.

A few recent microcontrollers beginning in 2003 have "self-programmable" flash memory[7].
Firmware

The earliest microcontrollers used hard-wired or mask ROM to store firmware. Later microcontrollers (such as the early versions of the Freescale 68HC11 and early PIC microcontrollers) had quartz windows that allowed ultraviolet light in to erase the EPROM.

The Microchip PIC16C84, introduced in 1993,[8] was the first microcontroller to use EEPROM to store firmware.

Also in 1993, Atmel introduced the first microcontroller using NOR Flash memory to store firmware.[7]

PSoC microcontrollers, introduced in 2002, store firmware in SONOS flash memory.

MRAM could potentially be used to store firmware.

Network switch

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Network switch

Typical SOHO network switch.

Back view of Atlantis network switch with Ethernet ports.

A network switch or switching hub is a computer networking device that connects network segments.
The term commonly refers to a network bridge that processes and routes data at the data link layer (layer 2) of the OSI model. Switches that additionally process data at the network layer (layer 3 and above) are often referred to as Layer 3 switches or multilayer switches.
The term network switch does not generally encompass unintelligent or passive network devices such as hubs and repeaters.
The first Ethernet switch was introduced by Kalpana in 1990.^[1]

Function

The network switch, packet switch (or just switch) plays an integral part in most Ethernet local area networks or LANs. Mid-to-large sized LANs contain a number of linked managed switches. Small office/home office (SOHO) applications typically use a single switch, or an all-purpose converged device such as gateway access to small office/home broadband services such as DSL router or cable Wi-Fi router. In most of these cases, the end user device contains a router and components that interface to the particular physical broadband technology, as in the Linksys 8-port and 48-port devices. User devices may also include a telephone interface to VoIP.
In the context of a standard 10/100 Ethernet switch, a switch operates at the data-link layer of the OSI model to create a different collision domain per switch port. If you have 4 computers A/B/C/D on 4 switch ports, then A and B can transfer data between them as well as C and D at the same time, and they will never interfere with each others' conversations. In the case of a "hub" then they would all have to share the bandwidth and run in Half duplex. The result is that there would be collisions and retransmissions. Using a switch is called micro-segmentation. It allows you to have dedicated bandwidth on point to point connections with every computer and to therefore run in Full duplex with no collisions.

Role of switches in networks

Switches may operate at one or more OSI layers, including physical, data link, network, or transport (i.e., end-to-end). A device that operates simultaneously at more than one of these layers is known as a multilayer switch.
In switches intended for commercial use, built-in or modular interfaces make it possible to connect different types of networks, including Ethernet, Fibre Channel, ATM, ITU-T G.hn and 802.11. This connectivity can be at any of the layers mentioned. While Layer 2 functionality is adequate for speed-shifting within one technology, interconnecting technologies such as Ethernet and token ring are easier at Layer 3.
Interconnection of different Layer 3 networks is done by routers. If there are any features that characterize "Layer-3 switches" as opposed to general-purpose routers, it tends to be that they are optimized, in larger switches, for high-density Ethernet connectivity.
In some service provider and other environments where there is a need for a great deal of analysis of network performance and security, switches may be connected between WAN routers as places for analytic modules. Some vendors provide firewall,^[2]^[3] network intrusion detection,^[4] and performance analysis modules that can plug into switch ports. Some of these functions may be on combined modules.^[5]
In other cases, the switch is used to create a mirror image of data that can go to an external device. Since most switch port mirroring provides only one mirrored stream, network hubs can be useful for fanning out data to several read-only analyzers, such as intrusion detection systems and packet sniffers.

Layer-specific functionality

Main article: Multilayer switch

A modular network switch with three network modules (a total of 24 Ethernet and 14 Fast Ethernet ports) and one power supply.

While switches may learn about topologies at many layers, and forward at one or more layers, they do tend to have common features. Other than for high-performance applications, modern commercial switches use primarily Ethernet interfaces, which can have different input and output speeds of 10, 100, 1000 or 10,000 megabits per second. Switch ports almost always default to Full duplex operation, unless there is a requirement for interoperability with devices that are strictly Half duplex. Half duplex means that the device can only send or receive at any given time, whereas Full duplex can send and receive at the same time.
At any layer, a modern switch may implement power over Ethernet (PoE), which avoids the need for attached devices, such as an IP telephone or wireless access point, to have a separate power supply. Since switches can have redundant power circuits connected to uninterruptible power supplies, the connected device can continue operating even when regular office power fails.

Layer-1 hubs versus higher-layer switches

A network hub, or repeater, is a fairly unsophisticated network device. Hubs do not manage any of the traffic that comes through them. Any packet entering a port is broadcast out or "repeated" on every other port, except for the port of entry. Since every packet is repeated on every other port, packet collisions result, which slows down the network.
There are specialized applications where a hub can be useful, such as copying traffic to multiple network sensors. High end switches have a feature which does the same thing called port mirroring. There is no longer any significant price difference between a hub and a low-end switch.^[6]

Layer 2

A network bridge, operating at the Media Access Control (MAC) sublayer of the data link layer, may interconnect a small number of devices in a home or office. This is a trivial case of bridging, in which the bridge learns the MAC address of each connected device. Single bridges also can provide extremely high performance in specialized applications such as storage area networks.
Classic bridges may also interconnect using a spanning tree protocol that disables links so that the resulting local area network is a tree without loops. In contrast to routers, spanning tree bridges must have topologies with only one active path between two points. The older IEEE 802.1D spanning tree protocol could be quite slow, with forwarding stopping for 30 seconds while the spanning tree would reconverge. A Rapid Spanning Tree Protocol was introduced as IEEE 802.1w, but the newest edition of IEEE 802.1D-2004, adopts the 802.1w extensions as the base standard. The IETF is specifying the TRILL protocol, which is the application of link-state routing technology to the layer-2 bridging problem. Devices which implement TRILL, called RBridges, combine the best features of both routers and bridges.
While "layer 2 switch" remains more of a marketing term than a technical term,^{[citation needed]} the products that were introduced as "switches" tended to use microsegmentation and Full duplex to prevent collisions among devices connected to Ethernets. By using an internal forwarding plane much faster than any interface, they give the impression of simultaneous paths among multiple devices.
Once a bridge learns the topology through a spanning tree protocol, it forwards data link layer frames using a layer 2 forwarding method. There are four forwarding methods a bridge can use, of which the second through fourth method were performance-increasing methods when used on "switch" products with the same input and output port speeds:

Store and forward: The switch buffers and, typically, performs a checksum on each frame before forwarding it on.
Cut through: The switch reads only up to the frame's hardware address before starting to forward it. There is no error checking with this method.
Fragment free: A method that attempts to retain the benefits of both "store and forward" and "cut through". Fragment free checks the first 64 bytes of the frame, where addressing information is stored. According to Ethernet specifications, collisions should be detected during the first 64 bytes of the frame, so frames that are in error because of a collision will not be forwarded. This way the frame will always reach its intended destination. Error checking of the actual data in the packet is left for the end device in Layer 3 or Layer 4 (OSI), typically a router.
Adaptive switching: A method of automatically switching between the other three modes.

Cut-through switches have to fall back to store and forward if the outgoing port is busy at the time the packet arrives. While there are specialized applications, such as storage area networks, where the input and output interfaces are the same speed, this is rarely the case in general LAN applications. In LANs, a switch used for end user access typically concentrates lower speed (e.g., 10/100 Mbit/s) into a higher speed (at least 1 Gbit/s). Alternatively, a switch that provides access to server ports usually connects to them at a much higher speed than is used by end user devices.Cypress Semiconductor design and manufacturing company along with TPACK offers the flexibility to cope with various system architecture for Ethernet switches through reference design. The reference design involves TPX4004 and CY7C15632KV18 72-Mbit SRAMs.

Layer 3

Within the confines of the Ethernet physical layer, a layer 3 switch can perform some or all of the functions normally performed by a router. A true router is able to forward traffic from one type of network connection (e.g., T1, DSL) to another (e.g., Ethernet, WiFi).
The most common layer-3 capability is awareness of IP multicast. With this awareness, a layer-3 switch can increase efficiency by delivering the traffic of a multicast group only to ports where the attached device has signaled that it wants to listen to that group. If a switch is not aware of multicasting and broadcasting, frames are also forwarded on all ports of each broadcast domain, but in the case of IP multicast this causes inefficient use of bandwidth. To work around this problem some switches implement IGMP snooping.^[7]

Layer 4

While the exact meaning of the term Layer-4 switch is vendor-dependent, it almost always starts with a capability for network address translation, but then adds some type of load distribution based on TCP sessions.^[8]
The device may include a stateful firewall, a VPN concentrator, or be an IPSec security gateway.

Layer 7

Layer 7 switches may distribute loads based on URL or by some installation-specific technique to recognize application-level transactions. A Layer-7 switch may include a web cache and participate in a content delivery network.^[9]

Rack-mounted 24-port 3Com switch

Types of switches

Form factor

Desktop, not mounted in an enclosure, typically intended to be used in a home or office environment outside of a wiring closet
Rack mounted
Chassis — with swappable "switch module" cards. e.g. Alcatel's OmniSwitch 7000; Cisco Catalyst switch 4500 and 6500; 3Com 7700, 7900E, 8800.
DIN rail mounted, normally seen in industrial environments or panels

Configuration options

Unmanaged switches — These switches have no configuration interface or options. They are plug and play. They are typically the least expensive switches, found in home, SOHO, or small businesses. They can be desktop or rack mounted.
Managed switches — These switches have one or more methods to modify the operation of the switch. Common management methods include: a serial console or command line interface accessed via telnet or Secure Shell, an embedded Simple Network Management Protocol (SNMP) agent allowing management from a remote console or management station, or a web interface for management from a web browser. Examples of configuration changes that one can do from a managed switch include: enable features such as Spanning Tree Protocol, set port speed, create or modify Virtual LANs (VLANs), etc. Two sub-classes of managed switches are marketed today:
- Smart (or intelligent) switches — These are managed switches with a limited set of management features. Likewise "web-managed" switches are switches which fall in a market niche between unmanaged and managed. For a price much lower than a fully managed switch they provide a web interface (and usually no CLI access) and allow configuration of basic settings, such as VLANs, port-speed and duplex.^[10]
- Enterprise Managed (or fully managed) switches — These have a full set of management features, including Command Line Interface, SNMP agent, and web interface. They may have additional features to manipulate configurations, such as the ability to display, modify, backup and restore configurations. Compared with smart switches, enterprise switches have more features that can be customized or optimized, and are generally more expensive than "smart" switches. Enterprise switches are typically found in networks with larger number of switches and connections, where centralized management is a significant savings in administrative time and effort. A stackable switch is a version of enterprise-managed switch.

Traffic monitoring on a switched network

Unless port mirroring or other methods such as RMON or SMON are implemented in a switch,^[11] it is difficult to monitor traffic that is bridged using a switch because all ports are isolated until one transmits data, and even then only the sending and receiving ports can see the traffic. These monitoring features rarely are present on consumer-grade switches.
Two popular methods that are specifically designed to allow a network analyst to monitor traffic are:

Port mirroring — the switch sends a copy of network packets to a monitoring network connection.
SMON — "Switch Monitoring" is described by RFC 2613 and is a protocol for controlling facilities such as port mirroring.

Another method to monitor may be to connect a Layer-1 hub between the monitored device and its switch port. This will induce minor delay, but will provide multiple interfaces that can be used to monitor the individual switch port.

Typical switch management features

Linksys 48-port switch

A rack-mounted switch with network cables

Turn some particular port range on or off
Link speed and duplex settings
Priority settings for ports
MAC filtering and other types of "port security" features which prevent MAC flooding
Use of Spanning Tree Protocol
SNMP monitoring of device and link health
Port mirroring (also known as: port monitoring, spanning port, SPAN port, roving analysis port or link mode port)
Link aggregation (also known as bonding, trunking or teaming)
VLAN settings
802.1X network access control
IGMP snooping

Link aggregation allows the use of multiple ports for the same connection achieving higher data transfer speeds. Creating VLANs can serve security and performance goals by reducing the size of the broadcast domain.

Network address translation ( NAT )

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Network address translation

In computer networking, network address translation (NAT) is the process of modifying network address information in datagram (IP) packet headers while in transit across a traffic routing device for the purpose of remapping a given address space into another.

Most often today, NAT is used in conjunction with network masquerading (or IP masquerading) which is a technique that hides an entire address space, usually consisting of private network addresses (RFC 1918), behind a single IP address in another, often public address space. This mechanism is implemented in a routing device that uses stateful translation tables to map the "hidden" addresses into a single address and then rewrites the outgoing Internet Protocol (IP) packets on exit so that they appear to originate from the router. In the reverse communications path, responses are mapped back to the originating IP address using the rules ("state") stored in the translation tables. The translation table rules established in this fashion are flushed after a short period without new traffic refreshing their state.

As described, the method enables communication through the router only when the conversation originates in the masqueraded network, since this establishes the translation tables. For example, a web browser in the masqueraded network can browse a website outside, but a web browser outside could not browse a web site in the masqueraded network. However, most NAT devices today allow the network administrator to configure translation table entries for permanent use. This feature is often referred to as "static NAT" or port forwarding and allows traffic originating in the 'outside' network to reach designated hosts in the masqueraded network.

Because of the popularity of this technique (see below), the term NAT has become virtually synonymous with the method of IP masquerading.

Network address translation has serious consequences (Drawbacks, Benefits) on the quality of Internet connectivity and requires careful attention to the details of its implementation. As a result, many methods have been devised to alleviate the issues encountered. See article on NAT traversal.

Overview

In the mid-1990s NAT became a popular tool for alleviating the IPv4 address exhaustion. It has become a standard, indispensable feature in routers for home and small-office Internet connections.

Most systems using NAT do so in order to enable multiple hosts on a private network to access the Internet using a single public IP address (see gateway). However, NAT breaks the originally envisioned model of IP end-to-end connectivity across the Internet, introduces complications in communication between hosts, and affects performance.

NAT obscures an internal network's structure: all traffic appears to outside parties as if it originated from the gateway machine.

Network address translation involves re-writing the source and/or destination IP addresses and usually also the TCP/UDP port numbers of IP packets as they pass through the NAT. Checksums (both IP and TCP/UDP) must also be rewritten to take account of the changes.

In a typical configuration, a local network uses one of the designated "private" IP address subnets (the RFC 1918). Private Network Addresses are 192.168.x.x, 172.16.x.x through 172.31.x.x, and 10.x.x.x (or using CIDR notation, 192.168/16, 172.16/12, and 10/8), and a router on that network has a private address (such as 192.168.0.1) in that address space. The router is also connected to the Internet with a single "public" address (known as "overloaded" NAT) or multiple "public" addresses assigned by an ISP. As traffic passes from the local network to the Internet, the source address in each packet is translated on the fly from the private addresses to the public address(es). The router tracks basic data about each active connection (particularly the destination address and port). When a reply returns to the router, it uses the connection tracking data it stored during the outbound phase to determine where on the internal network to forward the reply; the TCP or UDP client port numbers are used to demultiplex the packets in the case of overloaded NAT, or IP address and port number when multiple public addresses are available, on packet return. To a system on the Internet, the router itself appears to be the source/destination for this traffic.

Basic NAT and PAT

There are two levels of network address translation.

Basic NAT. This involves IP address translation only, not port mapping.
PAT (Port Address Translation). Also called simply "NAT" or "Network Address Port Translation, NAPT". This involves the translation of both IP addresses and port numbers.

All Internet packets have a source IP address and a destination IP address. Both or either of the source and destination addresses may be translated.

Some Internet packets do not have port numbers. For example, ICMP packets have no port numbers. However, the vast bulk of Internet traffic is TCP and UDP packets, which do have port numbers. Packets which do have port numbers have both a source port number and a destination port number. Both or either of the source and destination ports may be translated.

NAT which involves translation of the source IP address and/or source port is called source NAT or SNAT. This re-writes the IP address and/or port number of the computer which originated the packet.

NAT which involves translation of the destination IP address and/or destination port number is called destination NAT or DNAT. This re-writes the IP address and/or port number corresponding to the destination computer.

SNAT and DNAT may be applied simultaneously to Internet packets.

[edit] Types of NAT

Network address translation is implemented in a variety of schemes of translating addresses and port numbers, each affecting application communication protocols differently. In some application protocols that use IP address information, the application running on a node in the masqueraded network needs to determine the external address of the NAT, i.e., the address that its communication peers detect, and, furthermore, often needs to examine and categorize the type of mapping in use. For this purpose, the Simple traversal of UDP over NATs (STUN) protocol was developed (RFC 3489, March 2003). It classified NAT implementation as full cone NAT, (address) restricted cone NAT, port restricted cone NAT or symmetric NAT^[1] and proposed a methodology for testing a device accordingly. However, these procedures have since been deprecated from standards status, as the methods have proven faulty and inadequate to correctly assess many devices. New methods have been standardized in RFC 5389 (October 2008) and the STUN acronym now represents the new title of the specification: Session Traversal Utilities for NAT.

Full cone NAT, also known as one-to-one NAT Once an internal address (iAddr:iPort) is mapped to an external address (eAddr:ePort), any packets from iAddr:iPort will be sent through eAddr:ePort. Any external host can send packets to iAddr:iPort by sending packets to eAddr:ePort.
(Address) Restricted cone NAT Once an internal address (iAddr:iPort) is mapped to an external address (eAddr:ePort), any packets from iAddr:iPort will be sent through eAddr:ePort. An external host (hAddr:any) can send packets to iAddr:iPort by sending packets to eAddr:ePort only if iAddr:iPort had previously sent a packet to hAddr:any. "any" means the port number doesn't matter.
Port-Restricted cone NAT Like an (Address) Restricted cone NAT, but the restriction includes port numbers. Once an internal address (iAddr:iPort) is mapped to an external address (eAddr:ePort), any packets from iAddr:iPort will be sent through eAddr:ePort. An external host (hAddr:hPort) can send packets to iAddr:iPort by sending packets to eAddr:ePort only if iAddr:iPort had previously sent a packet to hAddr:hPort.
Symmetric NAT Each request from the same internal IP address and port to a specific destination IP address and port is mapped to a unique external source IP address and port. If the same internal host sends a packet even with the same source address and port but to a different destination, a different mapping is used. Only an external host that receives a packet from an internal host can send a packet back.

This terminology has been the source of much confusion, as it has proven inadequate at describing real-life NAT behavior.^[2] Many NAT implementations combine these types, and it is therefore better to refer to specific individual NAT behaviors instead of using the Cone/Symmetric terminology. Especially, most NAT translators combine symmetric NAT for outgoing connections with static port mapping, where incoming packets to the external address and port are redirected to a specific internal address and port. Some products can redirect packets to several internal hosts, e.g. to divide the load between a few servers. However, this introduces problems with more sophisticated communications that have many interconnected packets, and thus is rarely used.

Many NAT implementations follow the port preservation design. For most communications, they use the same values as internal and external port numbers. However, if two internal hosts attempt to communicate with the same external host using the same port number, the external port number used by the second host will be chosen at random. Such NAT will be sometimes perceived as (address) restricted cone NAT and other times as symmetric NAT.

NAT and TCP/UDP

"Pure NAT", operating on IP alone, may or may not correctly parse protocols that are totally concerned with IP information, such as ICMP, depending on whether the payload is interpreted by a host on the "inside" or "outside" of translation. As soon as the protocol stack is climbed, even with such basic protocols as TCP and UDP, the protocols will break unless NAT takes action beyond the network layer.

IP has a checksum in each packet header, which provides error detection only for the header. IP datagrams may become fragmented and it is necessary for a NAT to reassemble these fragments to allow correct recalculation of higher level checksums and correct tracking of which packets belong to which connection.

The major transport layer protocols, TCP and UDP, have a checksum that covers all the data they carry, as well as the TCP/UDP header, plus a "pseudo-header" that contains the source and destination IP addresses of the packet carrying the TCP/UDP header. For an originating NAT to successfully pass TCP or UDP, it must recompute the TCP/UDP header checksum based on the translated IP addresses, not the original ones, and put that checksum into the TCP/UDP header of the first packet of the fragmented set of packets. The receiving NAT must recompute the IP checksum on every packet it passes to the destination host, and also recognize and recompute the TCP/UDP header using the retranslated addresses and pseudo-header. This is not a completely solved problem. One solution is for the receiving NAT to reassemble the entire segment and then recompute a checksum calculated across all packets.

Originating host may perform Maximum transmission unit (MTU) path discovery (RFC 1191) to determine the packet size that can be transmitted without fragmentation, and then set the "don't fragment" bit in the appropriate packet header field.

Destination network address translation (DNAT)

DNAT is a technique for transparently changing the destination IP address of an en-route packet and performing the inverse function for any replies. Any router situated between two endpoints can perform this transformation of the packet.

DNAT is commonly used to publish a service located in a private network on a publicly accessible IP address. This use of DNAT is also called port forwarding.

SNAT

The usage of the term SNAT varies by vendor. Many vendors have proprietary definitions for SNAT. A common definition is Source NAT, the counterpart of Destination NAT (DNAT).

Microsoft uses the term for Secure NAT, in regards to the ISA Server extension discussed below. For Cisco Systems, SNAT means Stateful NAT.

The Internet Engineering Task Force (IETF) defines SNAT as Softwires Network Address Translation. This type of NAT is named after the Softwires working group that is charged with the standardization of discovery, control and encapsulation methods for connecting IPv4 networks across IPv6 networks and IPv6 networks across IPv4 networks.

Dynamic network address translation

Dynamic NAT, just like static NAT, is not common in smaller networks but is found within larger corporations with complex networks. The way dynamic NAT differs from static NAT is that where static NAT provides a one-to-one internal to public static IP address mapping, dynamic NAT doesn't make the mapping to the public IP address static and usually uses a group of available public IP addresses.

Applications affected by NAT

Some Application Layer protocols (such as FTP and SIP) send explicit network addresses within their application data. FTP in active mode, for example, uses separate connections for control traffic (commands) and for data traffic (file contents). When requesting a file transfer, the host making the request identifies the corresponding data connection by its network layer and transport layer addresses. If the host making the request lies behind a simple NAT firewall, the translation of the IP address and/or TCP port number makes the information received by the server invalid. The Session Initiation Protocol (SIP) controls Voice over IP (VoIP) communications and suffers the same problem. SIP may use multiple ports to set up a connection and transmit voice stream via RTP. IP addresses and port numbers are encoded in the payload data and must be known prior to the traversal of NATs. Without special techniques, such as STUN, NAT behavior is unpredictable and communications may fail.

Application Layer Gateway (ALG) software or hardware may correct these problems. An ALG software module running on a NAT firewall device updates any payload data made invalid by address translation. ALGs obviously need to understand the higher-layer protocol that they need to fix, and so each protocol with this problem requires a separate ALG.

Another possible solution to this problem is to use NAT traversal techniques using protocols such as STUN or ICE or proprietary approaches in a session border controller. NAT traversal is possible in both TCP- and UDP-based applications, but the UDP-based technique is simpler, more widely understood, and more compatible with legacy NATs. In either case, the high level protocol must be designed with NAT traversal in mind, and it does not work reliably across symmetric NATs or other poorly-behaved legacy NATs.

Other possibilities are UPnP (Universal Plug and Play) or Bonjour (NAT-PMP), but these require the cooperation of the NAT device.

Most traditional client-server protocols (FTP being the main exception), however, do not send layer 3 contact information and therefore do not require any special treatment by NATs. In fact, avoiding NAT complications is practically a requirement when designing new higher-layer protocols today.

NATs can also cause problems where IPsec encryption is applied and in cases where multiple devices such as SIP phones are located behind a NAT. Phones which encrypt their signaling with IPsec encapsulate the port information within the IPsec packet meaning that NA(P)T devices cannot access and translate the port. In these cases the NA(P)T devices revert to simple NAT operation. This means that all traffic returning to the NAT will be mapped onto one client causing the service to fail. There are a couple of solutions to this problem, one is to use TLS which operates at level 4 in the OSI Reference Model and therefore does not mask the port number, or to Encapsulate the IPsec within UDP - the latter being the solution chosen by TISPAN to achieve secure NAT traversal.

The DNS protocol vulnerability announced by Dan Kaminsky on 2008 July 8 is indirectly affected by NAT port mapping. To avoid DNS server cache poisoning, it is highly desirable to not translate UDP source port numbers of outgoing DNS requests from any DNS server which is behind a firewall which implements NAT. The recommended work-around for the DNS vulnerability is to make all caching DNS servers use randomized UDP source ports. If the NAT function de-randomizes the UDP source ports, the DNS server will be made vulnerable.

Drawbacks

Hosts behind NAT-enabled routers do not have end-to-end connectivity and cannot participate in some Internet protocols. Services that require the initiation of TCP connections from the outside network, or stateless protocols such as those using UDP, can be disrupted. Unless the NAT router makes a specific effort to support such protocols, incoming packets cannot reach their destination. Some protocols can accommodate one instance of NAT between participating hosts ("passive mode" FTP, for example), sometimes with the assistance of an application-level gateway (see below), but fail when both systems are separated from the Internet by NAT. Use of NAT also complicates tunneling protocols such as IPsec because NAT modifies values in the headers which interfere with the integrity checks done by IPsec and other tunneling protocols.

End-to-end connectivity has been a core principle of the Internet, supported for example by the Internet Architecture Board. Current Internet architectural documents observe that NAT is a violation of the End-to-End Principle, but that NAT does have a valid role in careful design.^[3] There is considerably more concern with the use of IPv6 NAT, and many IPv6 architects believe IPv6 was intended to remove the need for NAT.^[4]

Because of the short-lived nature of the stateful translation tables in NAT routers, devices on the internal network lose IP connectivity typically within a very short period of time unless they implement NAT keep-alive mechanisms by frequently accessing outside hosts. This dramatically shortens the power reserves on battery-operated hand-held devices and has thwarted more widespread deployment of such IP-native Internet-enabled devices.

Some Internet service providers (ISPs), especially in Russia, Asia and other "developing" regions provide their customers only with "local" IP addresses, due to limited number of external IP addresses allocated to those entities.^{[citation needed]} Thus, these customers must access services external to the ISP's network through NAT. As a result, the customers cannot achieve true end-to-end connectivity, in violation of the core principles of the Internet as laid out by the Internet Architecture Board.

Benefits

The primary benefit of IP-masquerading NAT is that it has been a practical solution to the impending exhaustion of IPv4 address space. Even large networks can be connected to the Internet with as little as a single IP address. The more common arrangement is having machines that require end-to-end connectivity supplied with a routable IP address, while having machines that do not provide services to outside users behind NAT with only a few IP addresses used to enable Internet access.

Some^[5] have also called this exact benefit a major drawback, since it delays the need for the implementation of IPv6, quote:

"... it is possible that its [NAT] widespread use will significantly delay the need to deploy IPv6. ... It is probably safe to say that networks would be better off without NAT, ..."

Examples of NAT software

IPFilter
PF (firewall): The OpenBSD Packet Filter.
Netfilter NAT engine
Internet Connection Sharing (ICS)
WinGate

Configuring Private VLANs

2010-05-14T22:23:00.000-07:00

Configuring Private VLANs

To configure a private VLAN, perform these steps:

Step 1 Set VTP mode to transparent.

Step 2 Create the primary and secondary VLANs and associate them. See the "Configuring and Associating VLANs in a Private VLAN" section.

Note If the VLAN is not created already, the private-VLAN configuration process creates it.

Step 3 Configure interfaces to be isolated or community host ports, and assign VLAN membership to the host port. See the "Configuring a Layer 2 Interface as a Private-VLAN Host Port" section.

Step 4 Configure interfaces as promiscuous ports, and map the promiscuous ports to the primary-secondary VLAN pair. See the "Configuring a Layer 2 Interface as a Private-VLAN Promiscuous Port" section.

Step 5 If inter-VLAN routing will be used, configure the primary SVI, and map secondary VLANs to the primary. See the "Mapping Secondary VLANs to a Primary VLAN Layer 3 VLAN Interface" section.

Step 6 Verify private-VLAN configuration.

Default Private-VLAN Configuration

No private VLANs are configured.

Private-VLAN Configuration Guidelines

Guidelines for configuring private VLANs fall into these categories:

•Secondary and Primary VLAN Configuration

•Private-VLAN Port Configuration

•Limitations with Other Features

Secondary and Primary VLAN Configuration

Follow these guidelines when configuring private VLANs:

•Set VTP to transparent mode. After you configure a private VLAN, you should not change the VTP mode to client or server. For information about VTP, see Chapter 14, "Configuring VTP."

•You must use VLAN configuration (config-vlan) mode to configure private VLANs. You cannot configure private VLANs in VLAN database configuration mode. For more information about VLAN configuration, see "VLAN Configuration Mode Options" section on page 13-7.

•After you have configured private VLANs, use the copy running-config startup config privileged EXEC command to save the VTP transparent mode configuration and private-VLAN configuration in the switch startup configuration file. Otherwise, if the switch resets, it defaults to VTP server mode, which does not support private VLANs.

•VTP does not propagate private-VLAN configuration. You must configure private VLANs on each device where you want private-VLAN ports.

•You cannot configure VLAN 1 or VLANs 1002 to 1005 as primary or secondary VLANs. Extended VLANs (VLAN IDs 1006 to 4094) can belong to private VLANs

•A primary VLAN can have one isolated VLAN and multiple community VLANs associated with it. An isolated or community VLAN can have only one primary VLAN associated with it.

•Although a private VLAN contains more than one VLAN, only one Spanning Tree Protocol (STP) instance runs for the entire private VLAN. When a secondary VLAN is associated with the primary VLAN, the STP parameters of the primary VLAN are propagated to the secondary VLAN.

•You can enable DHCP snooping on private VLANs. When you enable DHCP snooping on the primary VLAN, it is propagated to the secondary VLANs. If you configure DHCP on a secondary VLAN, the configuration does not take effect if the primary VLAN is already configured.

•When you enable IP source guard on private-VLAN ports, you must enable DHCP snooping on the primary VLAN.

•We recommend that you prune the private VLANs from the trunks on devices that carry no traffic in the private VLANs.

•You can apply different quality of service (QoS) configurations to primary, isolated, and community VLANs.

•When you configure private VLANs, sticky Address Resolution Protocol (ARP) is enabled by default, and ARP entries learned on Layer 3 private VLAN interfaces are sticky ARP entries. For security reasons, private VLAN port sticky ARP entries do not age out.

Note We recommend that you display and verify private-VLAN interface ARP entries.

Connecting a device with a different MAC address but with the same IP address generates a message and the ARP entry is not created. Because the private-VLAN port sticky ARP entries do not age out, you must manually remove private-VLAN port ARP entries if a MAC address changes.

–You can remove a private-VLAN ARP entry by using the no arp ip-address global configuration command.

–You can add a private-VLAN ARP entry by using the arp ip-address hardware-address type global configuration command.

•You can configure VLAN maps on primary and secondary VLANs (see the "Configuring VLAN Maps" section). However, we recommend that you configure the same VLAN maps on private-VLAN primary and secondary VLANs.

•When a frame is Layer-2 forwarded within a private VLAN, the same VLAN map is applied at the ingress side and at the egress side. When a frame is routed from inside a private VLAN to an external port, the private-VLAN map is applied at the ingress side.

–For frames going upstream from a host port to a promiscuous port, the VLAN map configured on the secondary VLAN is applied.

–For frames going downstream from a promiscuous port to a host port, the VLAN map configured on the primary VLAN is applied.

To filter out specific IP traffic for a private VLAN, you should apply the VLAN map to both the primary and secondary VLANs.

•You can apply router ACLs only on the primary-VLAN SVIs. The ACL is applied to both primary and secondary VLAN Layer 3 traffic.

•Although private VLANs provide host isolation at Layer 2, hosts can communicate with each other at Layer 3.

•Private VLANs support these Switched Port Analyzer (SPAN) features:

–You can configure a private-VLAN port as a SPAN source port.

–You can use VLAN-based SPAN (VSPAN) on primary, isolated, and community VLANs or use SPAN on only one VLAN to separately monitor egress or ingress traffic.

Private-VLAN Port Configuration

Follow these guidelines when configuring private-VLAN ports:

•Use only the private-VLAN configuration commands to assign ports to primary, isolated, or community VLANs. Layer 2 access ports assigned to the VLANs that you configure as primary, isolated, or community VLANs are inactive while the VLAN is part of the private-VLAN configuration. Layer 2 trunk interfaces remain in the STP forwarding state.

•Do not configure ports that belong to a PAgP or LACP EtherChannel as private-VLAN ports. While a port is part of the private-VLAN configuration, any EtherChannel configuration for it is inactive.

•Enable Port Fast and BPDU guard on isolated and community host ports to prevent STP loops due to misconfigurations and to speed up STP convergence (see Chapter 20, "Configuring Optional Spanning-Tree Features"). When enabled, STP applies the BPDU guard feature to all Port Fast-configured Layer 2 LAN ports. Do not enable Port Fast and BPDU guard on promiscuous ports.

•If you delete a VLAN used in the private-VLAN configuration, the private-VLAN ports associated with the VLAN become inactive.

•Private-VLAN ports can be on different network devices if the devices are trunk-connected and the primary and secondary VLANs have not been removed from the trunk.

Limitations with Other Features

When configuring private VLANs, remember these limitations with other features:

Note In some cases, the configuration is accepted with no error messages, but the commands have no effect.

•Do not configure fallback bridging on switches with private VLANs.

•When IGMP snooping is enabled on the switch (the default), the switch stack supports no more than 20 private-VLAN domains.

•Do not configure a remote SPAN (RSPAN) VLAN as a private-VLAN primary or secondary VLAN.

For more information about SPAN, see Chapter 28, "Configuring SPAN and RSPAN."

•Do not configure private-VLAN ports on interfaces configured for these other features:

–dynamic-access port VLAN membership

–Dynamic Trunking Protocol (DTP)

–Port Aggregation Protocol (PAgP)

–Link Aggregation Control Protocol (LACP)

–Multicast VLAN Registration (MVR)

–voice VLAN

•A private-VLAN port cannot be a secure port and should not be configured as a protected port.

•You can configure IEEE 802.1x port-based authentication on a private-VLAN port, but do not configure 802.1x with port security, voice VLAN, or per-user ACL on private-VLAN ports.

•A private-VLAN host or promiscuous port cannot be a SPAN destination port. If you configure a SPAN destination port as a private-VLAN port, the port becomes inactive.

•If you configure a static MAC address on a promiscuous port in the primary VLAN, you must add the same static address to all associated secondary VLANs. If you configure a static MAC address on a host port in a secondary VLAN, you must add the same static MAC address to the associated primary VLAN. When you delete a static MAC address from a private-VLAN port, you must remove all instances of the configured MAC address from the private VLAN.

Note Dynamic MAC addresses learned in one VLAN of a private VLAN are replicated in the associated VLANs. For example, a MAC address learned in a secondary VLAN is replicated in the primary VLAN. When the original dynamic MAC address is deleted or aged out, the replicated addresses are removed from the MAC address table.

•Configure Layer 3 VLAN interfaces (SVIs) only for primary VLANs.

Configuring and Associating VLANs in a Private VLAN

Beginning in privileged EXEC mode, follow these steps to configure a private VLAN:

Note The private-vlan commands do not take effect until you exit VLAN configuration mode.


	Command	Purpose
Step 1	configure terminal	Enter global configuration mode.
Step 2	vtp mode transparent	Set VTP mode to transparent (disable VTP).
Step 3	vlan vlan-id	Enter VLAN configuration mode and designate or create a VLAN that will be the primary VLAN. The VLAN ID range is 2 to 1001 and 1006 to 4094.
Step 4	private-vlan primary	Designate the VLAN as the primary VLAN.
Step 5	exit	Return to global configuration mode.
Step 6	vlan vlan-id	(Optional) Enter VLAN configuration mode and designate or create a VLAN that will be an isolated VLAN. The VLAN ID range is 2 to 1001 and 1006 to 4094.
Step 7	private-vlan isolated	Designate the VLAN as an isolated VLAN.
Step 8	exit	Return to global configuration mode.
Step 9	vlan vlan-id	(Optional) Enter VLAN configuration mode and designate or create a VLAN that will be a community VLAN. The VLAN ID range is 2 to 1001 and 1006 to 4094.
Step 10	private-vlan community	Designate the VLAN as a community VLAN.
Step 11	exit	Return to global configuration mode.
Step 12	vlan vlan-id	Enter VLAN configuration mode for the primary VLAN designated in Step 2.
Step 13	private-vlan association [add \| remove] secondary_vlan_list	Associate the secondary VLANs with the primary VLAN.
Step 14	end	Return to privileged EXEC mode.
Step 15	show vlan private-vlan [type] or show interfaces status	Verify the configuration.
Step 16	copy running-config startup config	Save your entries in the switch startup configuration file. To save the private-VLAN configuration, you need to save the VTP transparent mode configuration and private-VLAN configuration in the switch startup configuration file. Otherwise, if the switch resets, it defaults to VTP server mode, which does not support private VLANs.

When you associate secondary VLANs with a primary VLAN, note this syntax information:

•The secondary_vlan_list parameter cannot contain spaces. It can contain multiple comma-separated items. Each item can be a single private-VLAN ID or a hyphenated range of private-VLAN IDs.

•The secondary_vlan_list parameter can contain multiple community VLAN IDs but only one isolated VLAN ID.

•Enter a secondary_vlan_list, or use the add keyword with a secondary_vlan_list to associate secondary VLANs with a primary VLAN.

•Use the remove keyword with a secondary_vlan_list to clear the association between secondary VLANs and a primary VLAN.

•The command does not take effect until you exit VLAN configuration mode.

This example shows how to configure VLAN 20 as a primary VLAN, VLAN 501 as an isolated VLAN, and VLANs 502 and 503 as community VLANs, to associate them in a private VLAN, and to verify the configuration:

Switch# configure terminal

Switch(config)# vlan 20

Switch(config-vlan)# private-vlan primary

Switch(config-vlan)# exit

Switch(config)# vlan 501

Switch(config-vlan)# private-vlan isolated

Switch(config-vlan)# exit

Switch(config)# vlan 502

Switch(config-vlan)# private-vlan community

Switch(config-vlan)# exit

Switch(config)# vlan 503

Switch(config-vlan)# private-vlan community

Switch(config-vlan)# exit

Switch(config)# vlan 20

Switch(config-vlan)# private-vlan association 501-503

Switch(config-vlan)# end

Switch(config)# show vlan private vlan

Primary Secondary Type              Ports

------- --------- ----------------- ------------------------------------------

20      501       isolated

20      502       community

20      503       community

20      504       non-operational

Configuring a Layer 2 Interface as a Private-VLAN Host Port

Beginning in privileged EXEC mode, follow these steps to configure a Layer 2 interface as a private-VLAN host port and to associate it with primary and secondary VLANs:

Note Isolated and community VLANs are both secondary VLANs.


	Command	Purpose
Step 1	configure terminal	Enter global configuration mode.
Step 2	interface interface-id	Enter interface configuration mode for the Layer 2 interface to be configured.
Step 3	switchport mode private-vlan host	Configure the Layer 2 port as a private-VLAN host port.
Step 4	switchport private-vlan host-association primary_vlan_id secondary_vlan_id	Associate the Layer 2 port with a private VLAN.
Step 5	end	Return to privileged EXEC mode.
Step 6	show interfaces [interface-id] switchport	Verify the configuration.
Step 7	copy running-config startup config	(Optional) Save your entries in the switch startup configuration file.

This example shows how to configure an interface as a private-VLAN host port, associate it with a private-VLAN pair, and verify the configuration:

Switch# configure terminal

Switch(config)# interface fastethernet1/0/22

Switch(config-if)# switchport mode private-vlan host

Switch(config-if)# switchport private-vlan host-association 20 25

Switch(config-if)# end

Switch# show interfaces fastethernet1/0/22 switchport

Name: Fa1/0/22

Switchport: Enabled

Administrative Mode: private-vlan host

Operational Mode: private-vlan host

Administrative Trunking Encapsulation: negotiate

Operational Trunking Encapsulation: native

Negotiation of Trunking: Off

Access Mode VLAN: 1 (default)

Trunking Native Mode VLAN: 1 (default)

Administrative Native VLAN tagging: enabled

Voice VLAN: none

Administrative private-vlan host-association: 20 (VLAN0020) 25 (VLAN0025)

Administrative private-vlan mapping: none

Administrative private-vlan trunk native VLAN: none

Administrative private-vlan trunk Native VLAN tagging: enabled

Administrative private-vlan trunk encapsulation: dot1q

Administrative private-vlan trunk normal VLANs: none

Administrative private-vlan trunk private VLANs: none

Operational private-vlan:

20 (VLAN0020) 25 (VLAN0025)

Configuring a Layer 2 Interface as a Private-VLAN Promiscuous Port

Beginning in privileged EXEC mode, follow these steps to configure a Layer 2 interface as a private-VLAN promiscuous port and map it to primary and secondary VLANs:

Note Isolated and community VLANs are both secondary VLANs.


	Command	Purpose
Step 1	configure terminal	Enter global configuration mode.
Step 2	interface interface-id	Enter interface configuration mode for the Layer 2 interface to be configured.
Step 3	switchport mode private-vlan promiscuous	Configure the Layer 2 port as a private-VLAN promiscuous port.
Step 4	switchport private-vlan mapping primary_vlan_id {add \| remove} secondary_vlan_list	Map the private-VLAN promiscuous port to a primary VLAN and to selected secondary VLANs.
Step 5	end	Return to privileged EXEC mode.
Step 6	show interfaces [interface-id] switchport	Verify the configuration.
Step 7	copy running-config startup config	(Optional) Save your entries in the switch startup configuration file.

When you configure a Layer 2 interface as a private-VLAN promiscuous port, note this syntax information:

•The secondary_vlan_list parameter cannot contain spaces. It can contain multiple comma-separated items. Each item can be a single private-VLAN ID or a hyphenated range of private-VLAN IDs.

•Enter a secondary_vlan_list, or use the add keyword with a secondary_vlan_list to map the secondary VLANs to the private-VLAN promiscuous port.

•Use the remove keyword with a secondary_vlan_list to clear the mapping between secondary VLANs and the private-VLAN promiscuous port.

This example shows how to configure an interface as a private-VLAN promiscuous port and map it to a private VLAN. The interface is a member of primary VLAN 20 and secondary VLANs 501 to 503 are mapped to it.

Switch# configure terminal

Switch(config)# interface fastethernet1/0/2

Switch(config-if)# switchport mode private-vlan promiscuous

Switch(config-if)# switchport private-vlan mapping 20 add 501-503

Switch(config-if)# end

Use the show vlan private-vlan or the show interface status privileged EXEC command to display primary and secondary VLANs and private-VLAN ports on the switch.

Mapping Secondary VLANs to a Primary VLAN Layer 3 VLAN Interface

If the private VLAN will be used for inter-VLAN routing, you configure an SVI for the primary VLAN and map secondary VLANs to the SVI.

Note Isolated and community VLANs are both secondary VLANs.

Beginning in privileged EXEC mode, follow these steps to map secondary VLANs to the SVI of a primary VLAN to allow Layer 3 switching of private-VLAN traffic:


	Command	Purpose
Step 1	configure terminal	Enter global configuration mode.
Step 2	interface vlan primary_vlan_id	Enter interface configuration mode for the primary VLAN, and configure the VLAN as an SVI. The VLAN ID range is 2 to 1001 and 1006 to 4094.
Step 3	private-vlan mapping [add \| remove] secondary_vlan_list	Map the secondary VLANs to the Layer 3 VLAN interface of a primary VLAN to allow Layer 3 switching of private-VLAN ingress traffic.
Step 4	end	Return to privileged EXEC mode.
Step 5	show interface private-vlan mapping	Verify the configuration.
Step 6	copy running-config startup config	(Optional) Save your entries in the switch startup configuration file.

Note The private-vlan mapping interface configuration command only affects private-VLAN traffic that is Layer 3 switched.

When you map secondary VLANs to the Layer 3 VLAN interface of a primary VLAN, note this syntax information:

•The secondary_vlan_list parameter cannot contain spaces. It can contain multiple comma-separated items. Each item can be a single private-VLAN ID or a hyphenated range of private-VLAN IDs.

•Enter a secondary_vlan_list, or use the add keyword with a secondary_vlan_list to map the secondary VLANs to the primary VLAN.

•Use the remove keyword with a secondary_vlan_list to clear the mapping between secondary VLANs and the primary VLAN.

This example shows how to map the interfaces of VLANs 501and 502 to primary VLAN 10, which permits routing of secondary VLAN ingress traffic from private VLANs 501 to 502:

Switch# configure terminal

Switch(config)# interface vlan 10

Switch(config-if)# private-vlan mapping 501-502

Switch(config-if)# end

Switch# show interfaces private-vlan mapping

Interface Secondary VLAN Type

--------- -------------- -----------------

vlan10    501            isolated

vlan10    502            community

Monitoring Private VLANs

Table 16-1 shows the privileged EXEC commands for monitoring private-VLAN activity.

Table 16-1 Private VLAN Monitoring Commands
Command	Purpose
show interfaces status	Displays the status of interfaces, including the VLANs to which they belongs.
show vlan private-vlan [type]	Display the private-VLAN information for the switch stack.
show interface switchport	Display private-VLAN configuration on interfaces.
show interface private-vlan mapping	Display information about the private-VLAN mapping for VLAN SVIs.

This is an example of the output from the show vlan private-vlan command:

Switch(config)# show vlan private-vlan

Primary Secondary Type              Ports

------- --------- ----------------- ------------------------------------------

10      501       isolated          Fa2/0/1, Gi3/0/1, Gi3/0/2

10      502       community         Fa2/0/11, Gi3/0/1, Gi3/0/4

10      503       non-operational